CVPR 2024 Papers

Object pose refinement is essential for robust object pose estimation. Previous work has made significant progress towards instance-level object pose refinement. Yet category-level pose refinement is a more challenging problem due to large shape variations within a category and the discrepancies between the target object and the shape prior. To address these challenges we introduce a novel architecture for category-level object pose refinement. Our approach integrates an HS-layer and learnable affine transformations which aims to enhance the extraction and alignment of geometric information. Additionally we introduce a cross-cloud transformation mechanism that efficiently merges diverse data sources. Finally we push the limits of our model by incorporating the shape prior information for translation and size error prediction. We conducted extensive experiments to demonstrate the effectiveness of the proposed framework. Through extensive quantitative experiments we demonstrate significant improvement over the baseline method by a large margin across all metrics.

Recent advancements in Large Vision-Language Models (VLMs) have shown great promise in natural image domains allowing users to hold a dialogue about given visual content. However such general-domain VLMs perform poorly for Remote Sensing (RS) scenarios leading to inaccurate or fabricated information when presented with RS domain-specific queries. Such a behavior emerges due to the unique challenges introduced by RS imagery. For example to handle high-resolution RS imagery with diverse scale changes across categories and many small objects region-level reasoning is necessary alongside holistic scene interpretation. Furthermore the lack of domain-specific multimodal instruction following data as well as strong backbone models for RS make it hard for the models to align their behavior with user queries. To address these limitations we propose GeoChat - the first versatile remote sensing VLM that offers multitask conversational capabilities with high-resolution RS images. Specifically GeoChat can not only answer image-level queries but also accepts region inputs to hold region-specific dialogue. Furthermore it can visually ground objects in its responses by referring to their spatial coordinates. To address the lack of domain-specific datasets we generate a novel RS multimodal instruction-following dataset by extending image-text pairs from existing diverse RS datasets. Leveraging this rich dataset we fine-tune our remote sensing VLM based on the LLaVA-1.5 architecture. We establish a comprehensive benchmark for RS multitask conversations and compare with a number of baseline methods. GeoChat demonstrates robust zero-shot performance on various remote sensing tasks e.g. image and region captioning visual question answering scene classification visually grounded conversations and referring object detection. Our codes will be open-sourced.

Point clouds captured by different sensors such as RGB-D cameras and LiDAR possess non-negligible domain gaps. Most existing methods design different network architectures and train separately on point clouds from various sensors. Typically point-based methods achieve outstanding performances on even-distributed dense point clouds from RGB-D cameras while voxel-based methods are more efficient for large-range sparse LiDAR point clouds. In this paper we propose geometry-to-voxel auxiliary learning to enable voxel representations to access point-level geometric information which supports better generalisation of the voxel-based backbone with additional interpretations of multi-sensor point clouds. Specifically we construct hierarchical geometry pools generated by a voxel-guided dynamic point network which efficiently provide auxiliary fine-grained geometric information adapted to different stages of voxel features. We conduct experiments on joint multi-sensor datasets to demonstrate the effectiveness of GeoAuxNet. Enjoying elaborate geometric information our method outperforms other models collectively trained on multi-sensor datasets and achieve competitive results with the-state-of-art experts on each single dataset.

Real-world vision tasks frequently suffer from the appearance of unexpected adverse weather conditions including rain haze snow and raindrops. In the last decade convolutional neural networks and vision transformers have yielded outstanding results in single-weather video removal. However due to the absence of appropriate adaptation most of them fail to generalize to other weather conditions. Although ViWS-Net is proposed to remove adverse weather conditions in videos with a single set of pre-trained weights it is seriously blinded by seen weather at train-time and degenerates when coming to unseen weather during test-time. In this work we introduce test-time adaptation into adverse weather removal in videos and propose the first framework that integrates test-time adaptation into the iterative diffusion reverse process. Specifically we devise a diffusion-based network with a novel temporal noise model to efficiently explore frame-correlated information in degraded video clips at training stage. During inference stage we introduce a proxy task named Diffusion Tubelet Self-Calibration to learn the primer distribution of test video stream and optimize the model by approximating the temporal noise model for online adaptation. Experimental results on benchmark datasets demonstrate that our Test-Time Adaptation method with Diffusion-based network(Diff-TTA) outperforms state-of-the-art methods in terms of restoring videos degraded by seen weather conditions. Its generalizable capability is validated with unseen weather conditions in synthesized and real-world videos.

We present GenesisTex a novel method for synthesizing textures for 3D geometries from text descriptions. GenesisTex adapts the pretrained image diffusion model to texture space by texture space sampling. Specifically we maintain a latent texture map for each viewpoint which is updated with predicted noise on the rendering of the corresponding viewpoint. The sampled latent texture maps are then decoded into a final texture map. During the sampling process we focus on both global and local consistency across multiple viewpoints: global consistency is achieved through the integration of style consistency mechanisms within the noise prediction network and low-level consistency is achieved by dynamically aligning latent textures. Finally we apply reference-based inpainting and img2img on denser views for texture refinement. Our approach overcomes the limitations of slow optimization in distillation-based methods and instability in inpainting-based methods. Experiments on meshes from various sources demonstrate that our method surpasses the baseline methods quantitatively and qualitatively.

Recent advances in generative models trained on large-scale datasets have made it possible to synthesize high-quality samples across various domains. Moreover the emergence of strong inversion networks enables not only a reconstruction of real-world images but also the modification of attributes through various editing methods. However in certain domains related to privacy issues e.g. human faces advanced generative models along with strong inversion methods can lead to potential misuses. In this paper we propose an essential yet under-explored task called generative identity unlearning which steers the model not to generate an image of a specific identity. In the generative identity unlearning we target the following objectives: (i) preventing the generation of images with a certain identity and (ii) preserving the overall quality of the generative model. To satisfy these goals we propose a novel framework Generative Unlearning for Any Identity (GUIDE) which prevents the reconstruction of a specific identity by unlearning the generator with only a single image. GUIDE consists of two parts: (i) finding a target point for optimization that un-identifies the source latent code and (ii) novel loss functions that facilitate the unlearning procedure while less affecting the learned distribution. Our extensive experiments demonstrate that our proposed method achieves state-of-the-art performance in the generative machine unlearning task. The code is available at https://github.com/KHU-AGI/GUIDE.

Traditional 3D content creation tools empower users to bring their imagination to life by giving them direct control over a scene's geometry appearance motion and camera path. Creating computer-generated videos however is a tedious manual process which can be automated by emerging text-to-video diffusion models. Despite great promise video diffusion models are difficult to control hindering users to apply their creativity rather than amplifying it. To address this challenge we present a novel approach that combines the controllability of dynamic 3D meshes with the expressivity and editability of emerging diffusion models. For this purpose our approach takes an animated low-fidelity rendered mesh as input and injects the ground truth correspondence information obtained from the dynamic mesh into various stages of a pre-trained text-to-image generation model to output high-quality and temporally consistent frames. We demonstrate our approach on various examples where motion can be obtained by animating rigged assets or changing the camera path.

In recent research significant attention has been devoted to the open-vocabulary object detection task aiming to generalize beyond the limited number of classes labeled during training and detect objects described by arbitrary category names at inference. Compared with conventional object detection open vocabulary object detection largely extends the object detection categories. However it relies on calculating the similarity between image regions and a set of arbitrary category names with a pretrained vision-and-language model. This implies that despite its open-set nature the task still needs the predefined object categories during the inference stage. This raises the question: What if we do not have exact knowledge of object categories during inference? In this paper we call such a new setting as generative open-ended object detection which is a more general and practical problem. To address it we formulate object detection as a generative problem and propose a simple framework named GenerateU which can detect dense objects and generate their names in a free-form way. Particularly we employ Deformable DETR as a region proposal generator with a language model translating visual regions to object names. To assess the free-form object detection task we introduce an evaluation method designed to quantitatively measure the performance of generative outcomes. Extensive experiments demonstrate strong zero-shot detection performance of our GenerateU. For example on the LVIS dataset our GenerateU achieves comparable results to the open-vocabulary object detection method GLIP even though the category names are not seen by GenerateU during inference. Code is available at: https://github.com/FoundationVision/GenerateU.

The astonishing development of single-photon cameras has created an unprecedented opportunity for scientific and industrial imaging. However the high data throughput generated by these 1-bit sensors creates a significant bottleneck for low-power applications. In this paper we explore the possibility of generating a color image from a single binary frame of a single-photon camera. We evidently find this problem being particularly difficult to standard colorization approaches due to the substantial degree of exposure variation. The core innovation of our paper is an exposure synthesis model framed under a neural ordinary differential equation (Neural ODE) that allows us to generate a continuum of exposures from a single observation. This innovation ensures consistent exposure in binary images that colorizers take on resulting in notably enhanced colorization. We demonstrate applications of the method in single-image and burst colorization and show superior generative performance over baselines. Project website can be found at https://vishal-s-p.github.io/projects/2023/generative_quanta_color.html

Social interaction is a fundamental aspect of human behavior and communication. The way individuals position themselves in relation to others also known as proxemics conveys social cues and affects the dynamics of social interaction. Reconstructing such interaction from images presents challenges because of mutual occlusion and the limited availability of large training datasets. To address this we present a novel approach that learns a prior over the 3D proxemics two people in close social interaction and demonstrate its use for single-view 3D reconstruction. We start by creating 3D training data of interacting people using image datasets with contact annotations. We then model the proxemics using a novel denoising diffusion model called BUDDI that learns the joint distribution over the poses of two people in close social interaction. Sampling from our generative proxemics model produces realistic 3D human interactions which we validate through a perceptual study. We use BUDDI in reconstructing two people in close proximity from an image without any contact annotation via an optimization approach that uses the diffusion model as a prior. Our approach recovers accurate 3D social interactions from noisy initial estimates outperforming state-of-the-art methods. Our code data and model are available at: muelea.github.io/buddi.

We present a method that uses a text-to-image model to generate consistent content across multiple image scales enabling extreme semantic zooms into a scene e.g. ranging from a wide-angle landscape view of a forest to a macro shot of an insect sitting on one of the tree branches. We achieve this through a joint multi-scale diffusion sampling approach that encourages consistency across different scales while preserving the integrity of each individual sampling process. Since each generated scale is guided by a different text prompt our method enables deeper levels of zoom than traditional super-resolution methods that may struggle to create new contextual structure at vastly different scales. We compare our method qualitatively with alternative techniques in image super-resolution and outpainting and show that our method is most effective at generating consistent multi-scale content.

Humans can easily solve multimodal tasks in context with only a few demonstrations or simple instructions which current multimodal systems largely struggle to imitate. In this work we demonstrate that by effectively scaling up generative multimodal models their task-agnostic in-context learning capabilities can be significantly enhanced. We introduce Emu2 a generative multimodal model with 37 billion parameters which serves as a base model and general-purpose interface for a variety of multimodal tasks. Emu2 not only achieves strong performance in few-shot setting but can also be instruct-tuned to follow specific instructions such as visual question answering and object-grounded image generation. Emu2 even emerges to solve tasks that require on-the-fly reasoning such as visual prompting which existing models are unlikely to handle. We identify additional tasks where Emu2's in-context learning can further improve and discuss its broader societal impact. Our code and models will be made publicly available to facilitate future research.

In class incremental learning (CIL) scenarios the phenomenon of catastrophic forgetting caused by the classifier's bias towards the current task has long posed a significant challenge. It is mainly caused by the characteristic of discriminative models. With the growing popularity of the generative multi-modal models we would explore replacing discriminative models with generative ones for CIL. However transitioning from discriminative to generative models requires addressing two key challenges. The primary challenge lies in transferring the generated textual information into the classification of distinct categories. Additionally it requires formulating the task of CIL within a generative framework. To this end we propose a novel generative multi-modal model (GMM) framework for class incremental learning. Our approach directly generates labels for images using an adapted generative model. After obtaining the detailed text we use a text encoder to extract text features and employ feature matching to determine the most similar label as the classification prediction. In the conventional CIL settings we achieve significantly better results in long-sequence task scenarios. Under the Few-shot CIL setting we have improved by at least 14% over the current state-of-the-art methods with significantly less forgetting.

Most existing image compression approaches perform transform coding in the pixel space to reduce its spatial redundancy. However they encounter difficulties in achieving both high-realism and high-fidelity at low bitrate as the pixel-space distortion may not align with human perception. To address this issue we introduce a Generative Latent Coding (GLC) architecture which performs transform coding in the latent space of a generative vector-quantized variational auto-encoder (VQ-VAE) instead of in the pixel space. The generative latent space is characterized by greater sparsity richer semantic and better alignment with human perception rendering it advantageous for achieving high-realism and high-fidelity compression. Additionally we introduce a categorical hyper module to reduce the bit cost of hyper-information and a code-prediction-based supervision to enhance the semantic consistency. Experiments demonstrate that our GLC maintains high visual quality with less than 0.04 bpp on natural images and less than 0.01 bpp on facial images. On the CLIC2020 test set we achieve the same FID as MS-ILLM with 45% fewer bits. Furthermore the powerful generative latent space enables various applications built on our GLC pipeline such as image restoration and style transfer.

We present an approach to modeling an image-space prior on scene motion. Our prior is learned from a collection of motion trajectories extracted from real video sequences depicting natural oscillatory dynamics of objects such as treesflowers candles and clothes swaying in the wind. We model dense long-term motion in the Fourier domain as spectral volumes which we find are well-suited to prediction with diffusion models. Given a single image our trained model uses a frequency-coordinated diffusion sampling process to predict a spectral volume which can be converted into a motion texture that spans an entire video. Along with an image-based rendering module the predicted motion representation can be used for a number of downstream applications such as turning still images into seamlessly looping videos or allowing users to interact with objects in real images producing realistic simulated dynamics (by interpreting the spectral volumes as image-space modal bases). See our project page for more results: generative-dynamics.github.io

Generative 3D part assembly involves understanding part relationships and predicting their 6-DoF poses for assembling a realistic 3D shape. Prior work often focus on the geometry of individual parts neglecting part-whole hierarchies of objects. Leveraging two key observations: 1) super-part poses provide strong hints about part poses and 2) predicting super-part poses is easier due to fewer super-parts we propose a part-whole-hierarchy message passing network for efficient 3D part assembly. We first introduce super-parts by grouping geometrically similar parts without any semantic labels. Then we employ a part-whole hierarchical encoder wherein a super-part encoder predicts latent super-part poses based on input parts. Subsequently we transform the point cloud using the latent poses feeding it to the part encoder for aggregating super-part information and reasoning about part relationships to predict all part poses. In training only ground-truth part poses are required. During inference the predicted latent poses of super-parts enhance interpretability. Experimental results on the PartNet dataset that our method achieves state-of-the-art performance in part and connectivity accuracy and enables an interpretable hierarchical part assembly.

This paper addresses the challenge of example-based non-stationary texture synthesis. We introduce a novel two-step approach wherein users first modify a reference texture using standard image editing tools yielding an initial rough target for the synthesis. Subsequently our proposed method termed "self-rectification" automatically refines this target into a coherent seamless texture while faithfully preserving the distinct visual characteristics of the reference exemplar. Our method leverages a pre-trained diffusion network and uses self-attention mechanisms to gradually align the synthesized texture with the reference ensuring the retention of the structures in the provided target. Through experimental validation our approach exhibits exceptional proficiency in handling non-stationary textures demonstrating significant advancements in texture synthesis when compared to existing state-of-the-art techniques. Code is available at https://github.com/xiaorongjun000/Self-Rectification

We introduce a new task of generating "Illustrated Instructions" i.e. visual instructions customized to a user's needs. We identify desiderata unique to this task and formalize it through a suite of automatic and human evaluation metrics designed to measure the validity consistency and efficacy of the generations. We combine the power of large language models (LLMs) together with strong text-to-image generation diffusion models to propose a simple approach called StackedDiffusion which generates such illustrated instructions given text as input. The resulting model strongly outperforms baseline approaches and state-of-the-art multimodal LLMs; and in 30% of cases users even prefer it to human-generated articles. Most notably it enables various new and exciting applications far beyond what static articles on the web can provide such as personalized instructions complete with intermediate steps and pictures in response to a user's individual situation.

Generating human motions from textual descriptions has gained growing research interest due to its wide range of applications. However only a few works consider human-scene interactions together with text conditions which is crucial for visual and physical realism. This paper focuses on the task of generating human motions in 3D indoor scenes given text descriptions of the human-scene interactions. This task presents challenges due to the multimodality nature of text scene and motion as well as the need for spatial reasoning. To address these challenges we propose a new approach that decomposes the complex problem into two more manageable sub-problems: (1) language grounding of the target object and (2) object-centric motion generation. For language grounding of the target object we leverage the power of large language models. For motion generation we design an object-centric scene representation for the generative model to focus on the target object thereby reducing the scene complexity and facilitating the modeling of the relationship between human motions and the object. Experiments demonstrate the better motion quality of our approach compared to baselines and validate our design choices. Code will be available at https://zju3dv.github.io/text_scene_motion.