CVPR 2024 论文列表

Lifting 2D diffusion for 3D generation is a challenging problem due to the lack of geometric prior and the complex entanglement of materials and lighting in natural images. Existing methods have shown promise by first creating the geometry through score-distillation sampling (SDS) applied to rendered surface normals followed by appearance modeling. However relying on a 2D RGB diffusion model to optimize surface normals is suboptimal due to the distribution discrepancy between natural images and normals maps leading to instability in optimization. In this paper recognizing that the normal and depth information effectively describe scene geometry and be automatically estimated from images we propose to learn a generalizable Normal-Depth diffusion model for 3D generation. We achieve this by training on the large-scale LAION dataset together with the generalizable image-to-depth and normal prior models. In an attempt to alleviate the mixed illumination effects in the generated materials we introduce an albedo diffusion model to impose data-driven constraints on the albedo component. Our experiments show that when integrated into existing text-to-3D pipelines our models significantly enhance the detail richness achieving state-of-the-art results. Our project page is at https://aigc3d.github.io/richdreamer/.

Recent Text-to-Image (T2I) generation models such as Stable Diffusion and Imagen have made significant progress in generating high-resolution images based on text descriptions. However many generated images still suffer from issues such as artifacts/implausibility misalignment with text descriptions and low aesthetic quality. Inspired by the success of Reinforcement Learning with Human Feedback (RLHF) for large language models prior works collected human-provided scores as feedback on generated images and trained a reward model to improve the T2I generation. In this paper we enrich the feedback signal by (i) marking image regions that are implausible or misaligned with the text and (ii) annotating which words in the text prompt are misrepresented or missing on the image. We collect such rich human feedback on 18K generated images (RichHF-18K) and train a multimodal transformer to predict the rich feedback automatically. We show that the predicted rich human feedback can be leveraged to improve image generation for example by selecting high-quality training data to finetune and improve the generative models or by creating masks with predicted heatmaps to inpaint the problematic regions. Notably the improvements generalize to models (Muse) beyond those used to generate the images on which human feedback data were collected (Stable Diffusion variants). The RichHF-18K data set will be released in our GitHub repository: https://github.com/google-research/google-research/tree/master/richhf_18k.

Pan-sharpening is a super-resolution problem that essentially relies on spectra fusion of panchromatic (PAN) images and low-resolution multi-spectral (LRMS) images. The previous methods have validated the effectiveness of information fusion in the Fourier space of the whole image. However they haven't fully explored the Fourier relationships at different hierarchies between PAN and LRMS images. To this end we propose a Hierarchical Frequency Integration Network (HFIN) to facilitate hierarchical Fourier information integration for pan-sharpening. Specifically our network consists of two designs: information stratification and information integration. For information stratification we hierarchically decompose PAN and LRMS information into spatial global Fourier and local Fourier information and fuse them independently. For information integration the above hierarchical fused information is processed to further enhance their relationships and undergo comprehensive integration. Our method extend a new space for exploring the relationships of PAN and LRMS images enhancing the integration of spatial-frequency information. Extensive experiments robustly validate the effectiveness of the proposed network showcasing its superior performance compared to other state-of-the-art methods and generalization in real-world scenes and other fusion tasks as a general image fusion framework. Code is available at https://github.com/JosephTiTan/HFIN.

This research focuses on the issue of single-image reflection removal (SIRR) in real-world conditions examining it from two angles: the collection pipeline of real reflection pairs and the perception of real reflection locations. We devise an advanced reflection collection pipeline that is highly adaptable to a wide range of real-world reflection scenarios and incurs reduced costs in collecting large-scale aligned reflection pairs. In the process we develop a large-scale high-quality reflection dataset named Reflection Removal in the Wild (RRW). RRW contains over 14950 high-resolution real-world reflection pairs a dataset forty-five times larger than its predecessors. Regarding perception of reflection locations we identify that numerous virtual reflection objects visible in reflection images are not present in the corresponding ground-truth images. This observation drawn from the aligned pairs leads us to conceive the Maximum Reflection Filter (MaxRF). The MaxRF could accurately and explicitly characterize reflection locations from pairs of images. Building upon this we design a reflection location-aware cascaded framework specifically tailored for SIRR. Powered by these innovative techniques our solution achieves superior performance than current leading methods across multiple real-world benchmarks. Codes and datasets are available at \href https://github.com/zhuyr97/Reflection_RemoVal_CVPR2024 \color blue here .

The field of image synthesis is currently flourishing due to the advancements in diffusion models. While diffusion models have been successful their computational intensity has prompted the pursuit of more efficient alternatives. As a representative work non-autoregressive Transformers (NATs) have been recognized for their rapid generation. However a major drawback of these models is their inferior performance compared to diffusion models. In this paper we aim to re-evaluate the full potential of NATs by revisiting the design of their training and inference strategies. Specifically we identify the complexities in properly configuring these strategies and indicate the possible sub-optimality in existing heuristic-driven designs. Recognizing this we propose to go beyond existing methods by directly solving the optimal strategies in an automatic framework. The resulting method named AutoNAT advances the performance boundaries of NATs notably and is able to perform comparably with the latest diffusion models with a significantly reduced inference cost. The effectiveness of AutoNAT is comprehensively validated on four benchmark datasets i.e. ImageNet-256 & 512 MS-COCO and CC3M. Code and pre-trained models will be available at https://github.com/LeapLabTHU/ImprovedNAT.

Content-aware graphic layout generation aims to automatically arrange visual elements along with a given content such as an e-commerce product image. In this paper we argue that the current layout generation approaches suffer from the limited training data for the high-dimensional layout structure. We show that a simple retrieval augmentation can significantly improve the generation quality. Our model which is named Retrieval-Augmented Layout Transformer (RALF) retrieves nearest neighbor layout examples based on an input image and feeds these results into an autoregressive generator. Our model can apply retrieval augmentation to various controllable generation tasks and yield high-quality layouts within a unified architecture. Our extensive experiments show that RALF successfully generates content-aware layouts in both constrained and unconstrained settings and significantly outperforms the baselines.

Classifier-Free Guidance (CFG) has been widely used in text-to-image diffusion models where the CFG scale is introduced to control the strength of text guidance on the whole image space. However we argue that a global CFG scale results in spatial inconsistency on varying semantic strengths and suboptimal image quality. To address this problem we present a novel approach Semantic-aware Classifier-Free Guidance (S-CFG) to customize the guidance degrees for different semantic units in text-to-image diffusion models. Specifically we first design a training-free semantic segmentation method to partition the latent image into relatively independent semantic regions at each denoising step. In particular the cross-attention map in the denoising U-net backbone is renormalized for assigning each patch to the corresponding token while the self-attention map is used to complete the semantic regions. Then to balance the amplification of diverse semantic units we adaptively adjust the CFG scales across different semantic regions to rescale the text guidance degrees into a uniform level. Finally extensive experiments demonstrate the superiority of S-CFG over the original CFG strategy on various text-to-image diffusion models without requiring any extra training cost. our codes are available at https://github.com/SmilesDZgk/S-CFG.

Vector-Quantized (VQ-based) generative models usually consist of two basic components i.e. VQ tokenizers and generative transformers. Prior research focuses on improving the reconstruction fidelity of VQ tokenizers but rarely examines how the improvement in reconstruction affects the generation ability of generative transformers. In this paper we find that improving the reconstruction fidelity of VQ tokenizers does not necessarily improve the generation. Instead learning to compress semantic features within VQ tokenizers significantly improves generative transformers' ability to capture textures and structures. We thus highlight two competing objectives of VQ tokenizers for image synthesis: semantic compression and details preservation. Different from previous work that prioritizes better details preservation we propose Semantic-Quantized GAN (SeQ-GAN) with two learning phases to balance the two objectives. In the first phase we propose a semantic-enhanced perceptual loss for better semantic compression. In the second phase we fix the encoder and codebook but finetune the decoder to achieve better details preservation. Our proposed SeQ-GAN significantly improves VQ-based generative models for both unconditional and conditional image generation. Specifically SeQ-GAN achieves a Frechet Inception Distance (FID) of 6.25 and Inception Score (IS) of 140.9 on 256x256 ImageNet generation a remarkable improvement over VIT-VQGAN which obtains 11.2 FID and 97.2 IS.

Few-shot segmentation remains challenging due to the limitations of its labeling information for unseen classes. Most previous approaches rely on extracting high-level feature maps from the frozen visual encoder to compute the pixel-wise similarity as a key prior guidance for the decoder. However such a prior representation suffers from coarse granularity and poor generalization to new classes since these high-level feature maps have obvious category bias. In this work we propose to replace the visual prior representation with the visual-text alignment capacity to capture more reliable guidance and enhance the model generalization. Specifically we design two kinds of training-free prior information generation strategy that attempts to utilize the semantic alignment capability of the Contrastive Language-Image Pre-training model (CLIP) to locate the target class. Besides to acquire more accurate prior guidance we build a high-order relationship of attention maps and utilize it to refine the initial prior information. Experiments on both the PASCAL-5i and COCO-20i datasets show that our method obtains a clearly substantial improvement and reaches the new state-of-the-art performance. The code is available on the project website.

The goal of interactive image segmentation is to delineate specific regions within an image via visual or language prompts. Low-latency and high-quality interactive segmentation with diverse prompts remain challenging for existing specialist and generalist models. Specialist models with their limited prompts and task-specific designs experience high latency because the image must be recomputed every time the prompt is updated due to the joint encoding of image and visual prompts. Generalist models exemplified by the Segment Anything Model (SAM) have recently excelled in prompt diversity and efficiency lifting image segmentation to the foundation model era. However for high-quality segmentations SAM still lags behind state-of-the-art specialist models despite SAM being trained with x100 more segmentation masks. In this work we delve deep into the architectural differences between the two types of models. We observe that dense representation and fusion of visual prompts are the key design choices contributing to the high segmentation quality of specialist models. In light of this we reintroduce this dense design into the generalist models to facilitate the development of generalist models with high segmentation quality. To densely represent diverse visual prompts we propose to use a dense map to capture five types: clicks boxes polygons scribbles and masks. Thus we propose SegNext a next-generation interactive segmentation approach offering low latency high quality and diverse prompt support. Our method outperforms current state-of-the-art methods on HQSeg-44K and DAVIS quantitatively and qualitatively.

As with many machine learning problems the progress of image generation methods hinges on good evaluation metrics. One of the most popular is the Frechet Inception Distance (FID). FID estimates the distance between a distribution of Inception-v3 features of real images and those of images generated by the algorithm. We highlight important drawbacks of FID: Inception's poor representation of the rich and varied content generated by modern text-to-image models incorrect normality assumptions and poor sample complexity. We call for a reevaluation of FID's use as the primary quality metric for generated images. We empirically demonstrate that FID contradicts human raters it does not reflect gradual improvement of iterative text-to-image models it does not capture distortion levels and that it produces inconsistent results when varying the sample size. We also propose an alternative new metric CMMD based on richer CLIP embeddings and the maximum mean discrepancy distance with the Gaussian RBF kernel. It is an unbiased estimator that does not make any assumptions on the probability distribution of the embeddings and is sample efficient. Through extensive experiments and analysis we demonstrate that FID-based evaluations of text-to-image models may be unreliable and that CMMD offers a more robust and reliable assessment of image quality.

The inherent generative power of denoising diffusion models makes them well-suited for image restoration tasks where the objective is to find the optimal high-quality image within the generative space that closely resembles the input image. We propose a method to adapt a pretrained diffusion model for image restoration by simply adding noise to the input image to be restored and then denoise. Our method is based on the observation that the space of a generative model needs to be constrained. We impose this constraint by finetuning the generative model with a set of anchor images that capture the characteristics of the input image. With the constrained space we can then leverage the sampling strategy used for generation to do image restoration. We evaluate against previous methods and show superior performances on multiple real-world restoration datasets in preserving identity and image quality. We also demonstrate an important and practical application on personalized restoration where we use a personal album as the anchor images to constrain the generative space. This approach allows us to produce results that accurately preserve high-frequency details which previous works are unable to do. Project webpage: https://gen2res.github.io.

We propose residual denoising diffusion models (RDDM) a novel dual diffusion process that decouples the traditional single denoising diffusion process into residual diffusion and noise diffusion. This dual diffusion framework expands the denoising-based diffusion models initially uninterpretable for image restoration into a unified and interpretable model for both image generation and restoration by introducing residuals. Specifically our residual diffusion represents directional diffusion from the target image to the degraded input image and explicitly guides the reverse generation process for image restoration while noise diffusion represents random perturbations in the diffusion process. The residual prioritizes certainty while the noise emphasizes diversity enabling RDDM to effectively unify tasks with varying certainty or diversity requirements such as image generation and restoration. We demonstrate that our sampling process is consistent with that of DDPM and DDIM through coefficient transformation and propose a partially path-independent generation process to better understand the reverse process. Notably our RDDM enables a generic UNet trained with only an L1 loss and a batch size of 1 to compete with state-of-the-art image restoration methods. We provide code and pre-trained models to encourage further exploration application and development of our innovative framework (https://github.com/nachifur/RDDM).

Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity from relatively modest CNNs to large Transformer architectures. Still monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout since their knowledge of the visual world is restricted by the data seen during training and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better more generalizable depth estimation. We introduce Marigold a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.

Significant advancements have been made in image editing with the recent advance of the Diffusion model. However most of the current methods primarily focus on global or subject-level modifications and often face limitations when it comes to editing specific objects when there are other objects coexisting in the scene given solely textual prompts. In response to this challenge we introduce an object-level generative task called Referring Image Editing (RIE) which enables the identification and editing of specific source objects in an image using text prompts. To tackle this task effectively we propose a tailored framework called ReferDiffusion. It aims to disentangle input prompts into multiple embeddings and employs a mixed-supervised multi-stage training strategy. To facilitate further research in this domain we introduce the RefCOCO-Edit dataset comprising images editing prompts source object segmentation masks and reference edited images for training and evaluation. Our extensive experiments demonstrate the effectiveness of our approach in identifying and editing target objects while conventional general image editing and region-based image editing methods have difficulties in this challenging task.

Image stitching from different captures often results in non-rectangular boundaries which is often considered unappealing. To solve non-rectangular boundaries current solutions involve cropping which discards image content inpainting which can introduce unrelated content or warping which can distort non-linear features and introduce artifacts. To overcome these issues we introduce a novel diffusion-based learning framework RecDiffusion for image stitching rectangling. This framework combines Motion Diffusion Models (MDM) to generate motion fields effectively transitioning from the stitched image's irregular borders to a geometrically corrected intermediary. Followed by Content Diffusion Models (CDM) for image detail refinement. Notably our sampling process utilizes a weighted map to identify regions needing correction during each iteration of CDM. Our RecDiffusion ensures geometric accuracy and overall visual appeal surpassing all previous methods in both quantitative and qualitative measures when evaluated on public benchmarks. Code is released at https://github.com/lhaippp/RecDiffusion.

Text-to-image customization which aims to synthesize text-driven images for the given subjects has recently revolutionized content creation. Existing works follow the pseudo-word paradigm i.e. represent the given subjects as pseudo-words and then compose them with the given text. However the inherent entangled influence scope of pseudo-words with the given text results in a dual-optimum paradox i.e. the similarity of the given subjects and the controllability of the given text could not be optimal simultaneously. We present RealCustom that for the first time disentangles similarity from controllability by precisely limiting subject influence to relevant parts only achieved by gradually narrowing real text word from its general connotation to the specific subject and using its cross-attention to distinguish relevance. Specifically RealCustom introduces a novel "train-inference" decoupled framework: (1) during training RealCustom learns general alignment between visual conditions to original textual conditions by a novel adaptive scoring module to adaptively modulate influence quantity; (2) during inference a novel adaptive mask guidance strategy is proposed to iteratively update the influence scope and influence quantity of the given subjects to gradually narrow the generation of the real text word. Comprehensive experiments demonstrate the superior real-time customization ability of RealCustom in the open domain achieving both unprecedented similarity of the given subjects and controllability of the given text for the first time. The project page is https://corleone-huang.github.io/realcustom/.

The recently increased role of mobile photography has raised the standards of on-device photo processing tremendously. Despite the latest advancements in camera hardware the mobile camera sensor area cannot be increased significantly due to physical constraints leading to a pixel size of 0.6--2.0 \mum which results in strong image noise even in moderate lighting conditions. In the era of deep learning one can train a CNN model to perform robust image denoising. However there is still a lack of a substantially diverse dataset for this task. To address this problem we introduce a novel Mobile Image Denoising Dataset (MIDD) comprising over 400000 noisy / noise-free image pairs captured under various conditions by 20 different mobile camera sensors. Additionally we propose a new DPreview test set consisting of data from 294 different cameras for precise model evaluation. Furthermore we present the efficient baseline model SplitterNet for the considered mobile image denoising task that achieves high numerical and visual results while being able to process 8MP photos directly on smartphone GPUs in under one second. Thereby outperforming models with similar runtimes. This model is also compatible with recent mobile NPUs demonstrating an even higher speed when deployed on them. The conducted experiments demonstrate high robustness of the proposed solution when applied to images from previously unseen sensors showing its high generalizability. The datasets code and models can be found on the official project website.

We present Readout Guidance a method for controlling text-to-image diffusion models with learned signals. Readout Guidance uses readout heads lightweight networks trained to extract signals from the features of a pre-trained frozen diffusion model at every timestep. These readouts can encode single-image properties such as pose depth and edges; or higher-order properties that relate multiple images such as correspondence and appearance similarity. Furthermore by comparing the readout estimates to a user-defined target and back-propagating the gradient through the readout head these estimates can be used to guide the sampling process. Compared to prior methods for conditional generation Readout Guidance requires significantly fewer added parameters and training samples and offers a convenient and simple recipe for reproducing different forms of conditional control under a single framework with a single architecture and sampling procedure. We showcase these benefits in the applications of drag-based manipulation identity-consistent generation and spatially aligned control.