CVPR 2024 论文列表

This paper proposes Retrieval-Enhanced Asymmetric Diffusion (READ) for image-based robot motion planning. Given an image of the scene READ retrieves an initial motion from a database of image-motion pairs and uses a diffusion model to refine the motion for the given scene. Unlike prior retrieval-based diffusion models that require long forward-reverse diffusion paths READ directly diffuses between the source (retrieved) and target motions resulting in an efficient diffusion path. A second contribution of READ is its use of asymmetric diffusion whereby it preserves the kinematic feasibility of the generated motion by forward diffusion in a low-dimensional latent space while achieving high-resolution motion by reverse diffusion in the original task space using cold diffusion. Experimental results on various manipulation tasks demonstrate that READ outperforms state-of-the-art planning methods while ablation studies elucidate the contributions of asymmetric diffusion.

Recent advancements in diffusion-based models have demonstrated significant success in generating images from text. However video editing models have not yet reached the same level of visual quality and user control. To address this we introduce RAVE a zero-shot video editing method that leverages pre-trained text-to-image diffusion models without additional training. RAVE takes an input video and a text prompt to produce high-quality videos while preserving the original motion and semantic structure. It employs a novel noise shuffling strategy leveraging spatio-temporal interactions between frames to produce temporally consistent videos faster than existing methods. It is also efficient in terms of memory requirements allowing it to handle longer videos. RAVE is capable of a wide range of edits from local attribute modifications to shape transformations. In order to demonstrate the versatility of RAVE we create a comprehensive video evaluation dataset ranging from object-focused scenes to complex human activities like dancing and typing and dynamic scenes featuring swimming fish and boats. Our qualitative and quantitative experiments highlight the effectiveness of RAVE in diverse video editing scenarios compared to existing methods. Our code dataset and videos can be found in \href https://rave-video-edit.github.io/.

Recently the authors of Zero-1-to-3 demonstrated that a latent diffusion model pretrained with Internet-scale data can not only address the single-view 3D object reconstruction task but can even attain SOTA results in it. However when applied to the task of single-view 3D clothed human reconstruction Zero-1-to-3 (and related models) are unable to compete with the corresponding SOTA methods in this field despite being trained on clothed human data. In this work we aim to tailor Zero-1-to-3's approach to the single-view 3D clothed human reconstruction task in a much more principled and structured manner. To this end we propose R-Cyclic Diffuser a framework that adapts Zero-1-to-3's novel approach to clothed human data by fusing it with a pixel-aligned implicit model. R-Cyclic Diffuser offers a total of three new contributions. The first and primary contribution is R-Cyclic Diffuser's cyclical conditioning mechanism for novel view synthesis. This mechanism directly addresses the view inconsistency problem faced by Zero-1-to-3 and related models. Secondly we further enhance this mechanism with two key features - Lateral Inversion Constraint and Cyclic Noise Selection. Both features are designed to regularize and restrict the randomness of outputs generated by a latent diffusion model. Thirdly we show how SMPL-X body priors can be incorporated in a latent diffusion model such that novel views of clothed human bodies can be generated much more accurately. Our experiments show that R-Cyclic Diffuser is able to outperform current SOTA methods in single-view 3D clothed human reconstruction both qualitatively and quantitatively. Our code is made publicly available at https://github.com/kcyt/r-cyclic-diffuser.

The quality of the prompts provided to text-to-image diffusion models determines how faithful the generated content is to the user's intent often requiring `prompt engineering'. To harness visual concepts from target images without prompt engineering current approaches largely rely on embedding inversion by optimizing and then mapping them to pseudo-tokens. However working with such high-dimensional vector representations is challenging because they lack semantics and interpretability and only allow simple vector operations when using them. Instead this work focuses on inverting the diffusion model to obtain interpretable language prompts directly. The challenge of doing this lies in the fact that the resulting optimization problem is fundamentally discrete and the space of prompts is exponentially large; this makes using standard optimization techniques such as stochastic gradient descent difficult. To this end we utilize a delayed projection scheme to optimize for prompts representative of the vocabulary space in the model. Further we leverage the findings that different timesteps of the diffusion process cater to different levels of detail in an image. The later noisy timesteps of the forward diffusion process correspond to the semantic information and therefore prompt inversion in this range provides tokens representative of the image semantics. We show that our approach can identify semantically interpretable and meaningful prompts for a target image which can be used to synthesize diverse images with similar content. We further illustrate the application of the optimized prompts in evolutionary image generation and concept removal.

Diffusion-based generative models have exhibited remarkable capability in the production of high-fidelity visual content such as images and videos. However their performance is significantly contingent upon the quality of textual inputs commonly referred to as "prompts". The process of traditional prompt engineering while effective necessitates empirical expertise and poses challenges for inexperienced users. In this paper we introduce PromptCoT an innovative enhancer that autonomously refines prompts for users. PromptCoT is designed based on the observation that prompts which resemble the textual information of high-quality images in the training set often lead to superior generation performance. Therefore we fine-tune the pre-trained Large Language Models (LLM) using a curated text dataset that solely comprises descriptions of high-quality visual content. By doing so the LLM can capture the distribution of high-quality training texts enabling it to generate aligned continuations and revisions to boost the original texts. Nonetheless one drawback of pre-trained LLMs is their tendency to generate extraneous or irrelevant information. We employ the Chain-of-Thought (CoT) mechanism to improve the alignment between the original text prompts and their refined versions. CoT can extract and amalgamate crucial information from the aligned continuation and revision enabling reasonable inferences based on the contextual cues to produce a more comprehensive and nuanced final output. Considering computational efficiency instead of allocating a dedicated LLM for prompt enhancement to each individual model or dataset we integrate adapters that facilitate dataset-specific adaptation leveraging a shared pre-trained LLM as the foundation for this process. With independent fine-tuning of these adapters we can adapt PromptCoT to new datasets while minimally increasing training costs and memory usage. We evaluate the effectiveness of PromptCoT by assessing its performance on widely-used latent diffusion models for image and video generation. The results demonstrate significant improvements in key performance metrics.

Text-to-image (T2I) research has grown explosively in the past year owing to the large-scale pre-trained diffusion models and many emerging personalization and editing approaches. Yet one pain point persists: the text prompt engineering and searching high-quality text prompts for customized results is more art than science. Moreover as commonly argued: "an image is worth a thousand words" - the attempt to describe a desired image with texts often ends up being ambiguous and cannot comprehensively cover delicate visual details hence necessitating more additional controls from the visual domain. In this paper we take a bold step forward: taking "Text" out of a pretrained T2I diffusion model to reduce the burdensome prompt engineering efforts for users. Our proposed framework Prompt-Free Diffusion relies on only visual inputs to generate new images: it takes a reference image as "context" an optional image structural conditioning and an initial noise with absolutely no text prompt. The core architecture behind the scene is Semantic Context Encoder (SeeCoder) substituting the commonly used CLIP-based or LLM-based text encoder. The reusability of SeeCoder also makes it a convenient drop-in component: one can also pre-train a SeeCoder in one T2I model and reuse it for another. Through extensive experiments Prompt-Free Diffusion is experimentally found to (i) outperform prior exemplar-based image synthesis approaches; (ii) perform on par with state-of-the-art T2I models using prompts following the best practice; and (iii) be naturally extensible to other downstream applications such as anime figure generation and virtual try-on with promising quality. Our code and models will be open-sourced.

Text-guided image editing finds applications in various creative and practical fields. While recent studies in image generation have advanced the field they often struggle with the dual challenges of coherent image transformation and context preservation. In response our work introduces prompt augmentation a method amplifying a single input prompt into several target prompts strengthening textual context and enabling localised image editing. Specifically we use the augmented prompts to delineate the intended manipulation area. We propose a Contrastive Loss tailored to driving effective image editing by displacing edited areas and drawing preserved regions closer. Acknowledging the continuous nature of image manipulations we further refine our approach by incorporating the similarity concept creating a Soft Contrastive Loss. The new losses are incorporated to the diffusion model demonstrating improved or competitive image editing results on public datasets and generated images over state-of-the-art approaches.

Generative AI (GenAI) is transforming creative workflows through the capability to synthesize and manipulate images via high-level prompts. Yet creatives are not well supported to receive recognition or reward for the use of their content in GenAI training. To this end we propose ProMark a causal attribution technique to attribute a synthetically generated image to its training data concepts like objects motifs templates artists or styles. The concept information is proactively embedded into the input training images using imperceptible watermarks and the diffusion models (unconditional or conditional) are trained to retain the corresponding watermarks in generated images. We show that we can embed as many as 2^ 16 unique watermarks into the training data and each training image can contain more than one watermark. ProMark can maintain image quality whilst outperforming correlation-based attribution. Finally several qualitative examples are presented providing the confidence that the presence of the watermark conveys a causative relationship between training data and synthetic images.

Diffusion models have achieved remarkable success in generating high-quality diverse and creative images. However in text-based image generation they often struggle to accurately capture the intended meaning of the text. For instance a specified object might not be generated or an adjective might incorrectly alter unintended objects. Moreover we found that relationships indicating possession between objects are frequently overlooked. Despite the diversity of users' intentions in text existing methods often focus on only some aspects of these intentions. In this paper we propose Predicated Diffusion a unified framework designed to more effectively express users' intentions. It represents the intended meaning as propositions using predicate logic and treats the pixels in attention maps as fuzzy predicates. This approach provides a differentiable loss function that offers guidance for the image generation process to better fulfill the propositions. Comparative evaluations with existing methods demonstrated that Predicated Diffusion excels in generating images faithful to various text prompts while maintaining high image quality as validated by human evaluators and pretrained image-text models.

We introduce Posterior Distillation Sampling (PDS) a novel optimization method for parametric image editing based on diffusion models. Existing optimization-based methods which leverage the powerful 2D prior of diffusion models to handle various parametric images have mainly focused on generation. Unlike generation editing requires a balance between conforming to the target attribute and preserving the identity of the source content. Recent 2D image editing methods have achieved this balance by leveraging the stochastic latent encoded in the generative process of diffusion models. To extend the editing capabilities of diffusion models shown in pixel space to parameter space we reformulate the 2D image editing method into an optimization form named PDS. PDS matches the stochastic latents of the source and the target enabling the sampling of targets in diverse parameter spaces that align with a desired attribute while maintaining the source's identity. We demonstrate that this optimization resembles running a generative process with the target attribute but aligning this process with the trajectory of the source's generative process. Extensive editing results in Neural Radiance Fields and Scalable Vector Graphics representations demonstrate that PDS is capable of sampling targets to fulfill the aforementioned balance across various parameter spaces.

Unsupervised landmarks discovery (ULD) for an object category is a challenging computer vision problem. In pursuit of developing a robust ULD framework we explore the potential of a recent paradigm of self-supervised learning algorithms known as diffusion models. Some recent works have shown that these models implicitly contain important correspondence cues. Towards harnessing the potential of diffusion models for ULD task we make the following core contributions. First we propose a ZeroShot ULD baseline based on simple clustering of random pixel locations with nearest neighbour matching. It delivers better results than the existing ULD methods. Second motivated by the ZeroShot performance we develop a ULD algorithm based on diffusion features using self-training and clustering which also outperforms prior methods by notable margins. Third we introduce a new proxy task based on generating latent pose codes and also propose a two-stage clustering mechanism to facilitate effective pseudo-labeling resulting in a significant performance improvement. Overall our approach consistently outperforms state-of-the-art methods on four challenging benchmarks AFLW MAFL CatHeads and LS3D by significant margins.

Recent advancements in personalized image generation using diffusion models have been noteworthy. However existing methods suffer from inefficiencies due to the requirement for subject-specific fine-tuning. This computationally intensive process hinders efficient deployment limiting practical usability. Moreover these methods often grapple with identity distortion and limited expression diversity. In light of these challenges we propose PortraitBooth an innovative approach designed for high efficiency robust identity preservation and expression-editable text-to-image generation without the need for fine-tuning. PortraitBooth leverages subject embeddings from a face recognition model for personalized image generation without fine-tuning. It eliminates computational overhead and mitigates identity distortion. The introduced dynamic identity preservation strategy further ensures close resemblance to the original image identity. Moreover PortraitBooth incorporates emotion-aware cross-attention control for diverse facial expressions in generated images supporting text-driven expression editing. Its scalability enables efficient and high-quality image creation including multi-subject generation. Extensive results demonstrate superior performance over other state-of-the-art methods in both single and multiple image generation scenarios.

We present PointInfinity an efficient family of point cloud diffusion models. Our core idea is to use a transformer-based architecture with a fixed-size resolution-invariant latent representation. This enables efficient training with low-resolution point clouds while allowing high-resolution point clouds to be generated during inference. More importantly we show that scaling the test-time resolution beyond the training resolution improves the fidelity of generated point clouds and surfaces. We analyze this phenomenon and draw a link to classifier-free guidance commonly used in diffusion models demonstrating that both allow trading off fidelity and variability during inference. Experiments on CO3D show that PointInfinity can efficiently generate high-resolution point clouds (up to 131k points 31 times more than Point-E) with state-of-the-art quality.

Pre-training a model and then fine-tuning it on downstream tasks has demonstrated significant success in the 2D image and NLP domains. However due to the unordered and non-uniform density characteristics of point clouds it is non-trivial to explore the prior knowledge of point clouds and pre-train a point cloud backbone. In this paper we propose a novel pre-training method called Point cloud Diffusion pre-training PointDif. We consider the point cloud pre-training task as a conditional point-to-point generation problem and introduce a conditional point generator. This generator aggregates the features extracted by the backbone and employs them as the condition to guide the point-to-point recovery from the noisy point cloud thereby assisting the backbone in capturing both local and global geometric priors as well as the global point density distribution of the object. We also present a recurrent uniform sampling optimization strategy which enables the model to uniformly recover from various noise levels and learn from balanced supervision. Our PointDif achieves substantial improvement across various real-world datasets for diverse downstream tasks such as classification segmentation and detection. Specifically PointDif attains 70.0% mIoU on S3DIS Area 5 for the segmentation task and achieves an average improvement of 2.4% on ScanObjectNN for the classification task compared to TAP. Furthermore our pre-training framework can be flexibly applied to diverse point cloud backbones and bring considerable gains. Code is available at https://github.com/zhengxiaozx/PointDif

Diffusion models have shown tremendous results in image generation. However due to the iterative nature of the diffusion process and its reliance on classifier-free guidance inference times are slow. In this paper we propose a new distillation approach for guided diffusion models in which an external lightweight guide model is trained while the original text-to-image model remains frozen. We show that our method reduces the inference computation of classifier-free guided latent-space diffusion models by almost half and only requires 1% trainable parameters of the base model. Furthermore once trained our guide model can be applied to various fine-tuned domain-specific versions of the base diffusion model without the need for additional training: this "plug-and-play" functionality drastically improves inference computation while maintaining the visual fidelity of generated images. Empirically we show that our approach is able to produce visually appealing results and achieve a comparable FID score to the teacher with as few as 8 to 16 steps.

With recent developments in Embodied Artificial Intelligence (EAI) research there has been a growing demand for high-quality large-scale interactive scene generation. While prior methods in scene synthesis have prioritized the naturalness and realism of the generated scenes the physical plausibility and interactivity of scenes have been largely left unexplored. To address this disparity we introduce PhyScene a novel method dedicated to generating interactive 3D scenes characterized by realistic layouts articulated objects and rich physical interactivity tailored for embodied agents. Based on a conditional diffusion model for capturing scene layouts we devise novel physics- and interactivity-based guidance mechanisms that integrate constraints from object collision room layout and object reachability. Through extensive experiments we demonstrate that PhyScene effectively leverages these guidance functions for physically interactable scene synthesis outperforming existing state-of-the-art scene synthesis methods by a large margin. Our findings suggest that the scenes generated by PhyScene hold considerable potential for facilitating diverse skill acquisition among agents within interactive environments thereby catalyzing further advancements in embodied AI research.

Diffusion models (DMs) embark a new era of generative modeling and offer more opportunities for efficient generating high-quality and realistic data samples. However their widespread use has also brought forth new challenges in model security which motivates the creation of more effective adversarial attackers on DMs to understand its vulnerability. We propose CAAT a simple but generic and efficient approach that does not require costly training to effectively fool latent diffusion models (LDMs). The approach is based on the observation that cross-attention layers exhibits higher sensitivity to gradient change allowing for leveraging subtle perturbations on published images to significantly corrupt the generated images. We show that a subtle perturbation on an image can significantly impact the cross-attention layers thus changing the mapping between text and image during the fine-tuning of customized diffusion models. Extensive experiments demonstrate that CAAT is compatible with diverse diffusion models and outperforms baseline attack methods in a more effective (more noise) and efficient (twice as fast as Anti-DreamBooth and Mist) manner.

We present personalized residuals and localized attention-guided sampling for efficient concept-driven generation using text-to-image diffusion models. Our method first represents concepts by freezing the weights of a pretrained text-conditioned diffusion model and learning low-rank residuals for a small subset of the model's layers. The residual-based approach then directly enables application of our proposed sampling technique which applies the learned residuals only in areas where the concept is localized via cross-attention and applies the original diffusion weights in all other regions. Localized sampling therefore combines the learned identity of the concept with the existing generative prior of the underlying diffusion model. We show that personalized residuals effectively capture the identity of a concept in 3 minutes on a single GPU without the use of regularization images and with fewer parameters than previous models and localized sampling allows using the original model as strong prior for large parts of the image.

Finding correspondences between 3D shapes is an important and long-standing problem in computer vision graphics and beyond. A prominent challenge are partial-to-partial shape matching settings which occur when the shapes to match are only observed incompletely (e.g. from 3D scanning). Although partial-to-partial matching is a highly relevant setting in practice it is rarely explored. Our work bridges the gap between existing (rather artificial) 3D full shape matching and partial-to-partial real-world settings by exploiting geometric consistency as a strong constraint. We demonstrate that it is indeed possible to solve this challenging problem in a variety of settings. For the first time we achieve geometric consistency for partial-to-partial matching which is realized by a novel integer non-linear program formalism building on triangle product spaces along with a new pruning algorithm based on linear integer programming. Further we generate a new inter-class dataset for partial-to-partial shape-matching. We show that our method outperforms current SOTA methods on both an established intra-class dataset and our novel inter-class dataset.