CVPR 2024 论文列表

In recent research significant attention has been devoted to the open-vocabulary object detection task aiming to generalize beyond the limited number of classes labeled during training and detect objects described by arbitrary category names at inference. Compared with conventional object detection open vocabulary object detection largely extends the object detection categories. However it relies on calculating the similarity between image regions and a set of arbitrary category names with a pretrained vision-and-language model. This implies that despite its open-set nature the task still needs the predefined object categories during the inference stage. This raises the question: What if we do not have exact knowledge of object categories during inference? In this paper we call such a new setting as generative open-ended object detection which is a more general and practical problem. To address it we formulate object detection as a generative problem and propose a simple framework named GenerateU which can detect dense objects and generate their names in a free-form way. Particularly we employ Deformable DETR as a region proposal generator with a language model translating visual regions to object names. To assess the free-form object detection task we introduce an evaluation method designed to quantitatively measure the performance of generative outcomes. Extensive experiments demonstrate strong zero-shot detection performance of our GenerateU. For example on the LVIS dataset our GenerateU achieves comparable results to the open-vocabulary object detection method GLIP even though the category names are not seen by GenerateU during inference. Code is available at: https://github.com/FoundationVision/GenerateU.

The astonishing development of single-photon cameras has created an unprecedented opportunity for scientific and industrial imaging. However the high data throughput generated by these 1-bit sensors creates a significant bottleneck for low-power applications. In this paper we explore the possibility of generating a color image from a single binary frame of a single-photon camera. We evidently find this problem being particularly difficult to standard colorization approaches due to the substantial degree of exposure variation. The core innovation of our paper is an exposure synthesis model framed under a neural ordinary differential equation (Neural ODE) that allows us to generate a continuum of exposures from a single observation. This innovation ensures consistent exposure in binary images that colorizers take on resulting in notably enhanced colorization. We demonstrate applications of the method in single-image and burst colorization and show superior generative performance over baselines. Project website can be found at https://vishal-s-p.github.io/projects/2023/generative_quanta_color.html

Social interaction is a fundamental aspect of human behavior and communication. The way individuals position themselves in relation to others also known as proxemics conveys social cues and affects the dynamics of social interaction. Reconstructing such interaction from images presents challenges because of mutual occlusion and the limited availability of large training datasets. To address this we present a novel approach that learns a prior over the 3D proxemics two people in close social interaction and demonstrate its use for single-view 3D reconstruction. We start by creating 3D training data of interacting people using image datasets with contact annotations. We then model the proxemics using a novel denoising diffusion model called BUDDI that learns the joint distribution over the poses of two people in close social interaction. Sampling from our generative proxemics model produces realistic 3D human interactions which we validate through a perceptual study. We use BUDDI in reconstructing two people in close proximity from an image without any contact annotation via an optimization approach that uses the diffusion model as a prior. Our approach recovers accurate 3D social interactions from noisy initial estimates outperforming state-of-the-art methods. Our code data and model are available at: muelea.github.io/buddi.

Humans can easily solve multimodal tasks in context with only a few demonstrations or simple instructions which current multimodal systems largely struggle to imitate. In this work we demonstrate that by effectively scaling up generative multimodal models their task-agnostic in-context learning capabilities can be significantly enhanced. We introduce Emu2 a generative multimodal model with 37 billion parameters which serves as a base model and general-purpose interface for a variety of multimodal tasks. Emu2 not only achieves strong performance in few-shot setting but can also be instruct-tuned to follow specific instructions such as visual question answering and object-grounded image generation. Emu2 even emerges to solve tasks that require on-the-fly reasoning such as visual prompting which existing models are unlikely to handle. We identify additional tasks where Emu2's in-context learning can further improve and discuss its broader societal impact. Our code and models will be made publicly available to facilitate future research.

In class incremental learning (CIL) scenarios the phenomenon of catastrophic forgetting caused by the classifier's bias towards the current task has long posed a significant challenge. It is mainly caused by the characteristic of discriminative models. With the growing popularity of the generative multi-modal models we would explore replacing discriminative models with generative ones for CIL. However transitioning from discriminative to generative models requires addressing two key challenges. The primary challenge lies in transferring the generated textual information into the classification of distinct categories. Additionally it requires formulating the task of CIL within a generative framework. To this end we propose a novel generative multi-modal model (GMM) framework for class incremental learning. Our approach directly generates labels for images using an adapted generative model. After obtaining the detailed text we use a text encoder to extract text features and employ feature matching to determine the most similar label as the classification prediction. In the conventional CIL settings we achieve significantly better results in long-sequence task scenarios. Under the Few-shot CIL setting we have improved by at least 14% over the current state-of-the-art methods with significantly less forgetting.

Most existing image compression approaches perform transform coding in the pixel space to reduce its spatial redundancy. However they encounter difficulties in achieving both high-realism and high-fidelity at low bitrate as the pixel-space distortion may not align with human perception. To address this issue we introduce a Generative Latent Coding (GLC) architecture which performs transform coding in the latent space of a generative vector-quantized variational auto-encoder (VQ-VAE) instead of in the pixel space. The generative latent space is characterized by greater sparsity richer semantic and better alignment with human perception rendering it advantageous for achieving high-realism and high-fidelity compression. Additionally we introduce a categorical hyper module to reduce the bit cost of hyper-information and a code-prediction-based supervision to enhance the semantic consistency. Experiments demonstrate that our GLC maintains high visual quality with less than 0.04 bpp on natural images and less than 0.01 bpp on facial images. On the CLIC2020 test set we achieve the same FID as MS-ILLM with 45% fewer bits. Furthermore the powerful generative latent space enables various applications built on our GLC pipeline such as image restoration and style transfer.

We present an approach to modeling an image-space prior on scene motion. Our prior is learned from a collection of motion trajectories extracted from real video sequences depicting natural oscillatory dynamics of objects such as treesflowers candles and clothes swaying in the wind. We model dense long-term motion in the Fourier domain as spectral volumes which we find are well-suited to prediction with diffusion models. Given a single image our trained model uses a frequency-coordinated diffusion sampling process to predict a spectral volume which can be converted into a motion texture that spans an entire video. Along with an image-based rendering module the predicted motion representation can be used for a number of downstream applications such as turning still images into seamlessly looping videos or allowing users to interact with objects in real images producing realistic simulated dynamics (by interpreting the spectral volumes as image-space modal bases). See our project page for more results: generative-dynamics.github.io

Generative 3D part assembly involves understanding part relationships and predicting their 6-DoF poses for assembling a realistic 3D shape. Prior work often focus on the geometry of individual parts neglecting part-whole hierarchies of objects. Leveraging two key observations: 1) super-part poses provide strong hints about part poses and 2) predicting super-part poses is easier due to fewer super-parts we propose a part-whole-hierarchy message passing network for efficient 3D part assembly. We first introduce super-parts by grouping geometrically similar parts without any semantic labels. Then we employ a part-whole hierarchical encoder wherein a super-part encoder predicts latent super-part poses based on input parts. Subsequently we transform the point cloud using the latent poses feeding it to the part encoder for aggregating super-part information and reasoning about part relationships to predict all part poses. In training only ground-truth part poses are required. During inference the predicted latent poses of super-parts enhance interpretability. Experimental results on the PartNet dataset that our method achieves state-of-the-art performance in part and connectivity accuracy and enables an interpretable hierarchical part assembly.

This paper addresses the challenge of example-based non-stationary texture synthesis. We introduce a novel two-step approach wherein users first modify a reference texture using standard image editing tools yielding an initial rough target for the synthesis. Subsequently our proposed method termed "self-rectification" automatically refines this target into a coherent seamless texture while faithfully preserving the distinct visual characteristics of the reference exemplar. Our method leverages a pre-trained diffusion network and uses self-attention mechanisms to gradually align the synthesized texture with the reference ensuring the retention of the structures in the provided target. Through experimental validation our approach exhibits exceptional proficiency in handling non-stationary textures demonstrating significant advancements in texture synthesis when compared to existing state-of-the-art techniques. Code is available at https://github.com/xiaorongjun000/Self-Rectification

We introduce a new task of generating "Illustrated Instructions" i.e. visual instructions customized to a user's needs. We identify desiderata unique to this task and formalize it through a suite of automatic and human evaluation metrics designed to measure the validity consistency and efficacy of the generations. We combine the power of large language models (LLMs) together with strong text-to-image generation diffusion models to propose a simple approach called StackedDiffusion which generates such illustrated instructions given text as input. The resulting model strongly outperforms baseline approaches and state-of-the-art multimodal LLMs; and in 30% of cases users even prefer it to human-generated articles. Most notably it enables various new and exciting applications far beyond what static articles on the web can provide such as personalized instructions complete with intermediate steps and pictures in response to a user's individual situation.

Generating human motions from textual descriptions has gained growing research interest due to its wide range of applications. However only a few works consider human-scene interactions together with text conditions which is crucial for visual and physical realism. This paper focuses on the task of generating human motions in 3D indoor scenes given text descriptions of the human-scene interactions. This task presents challenges due to the multimodality nature of text scene and motion as well as the need for spatial reasoning. To address these challenges we propose a new approach that decomposes the complex problem into two more manageable sub-problems: (1) language grounding of the target object and (2) object-centric motion generation. For language grounding of the target object we leverage the power of large language models. For motion generation we design an object-centric scene representation for the generative model to focus on the target object thereby reducing the scene complexity and facilitating the modeling of the relationship between human motions and the object. Experiments demonstrate the better motion quality of our approach compared to baselines and validate our design choices. Code will be available at https://zju3dv.github.io/text_scene_motion.

In this paper we explore a novel challenging generation task i.e. Handwritten Mathematical Expression Generation (HMEG) from symbolic sequences. Since symbolic sequences are naturally graph-structured data we formulate HMEG as a graph-to-image (G2I) generation problem. Unlike the generation of natural images HMEG requires critic layout clarity for synthesizing correct and recognizable formulas but has no real masks available to supervise the learning process. To alleviate this challenge we propose a novel end-to-end G2I generation pipeline (i.e. graph - layout - mask - image) which requires no real masks or nondifferentiable alignment between layouts and masks. Technically to boost the capacity of predicting detailed relations among adjacent symbols we propose a Less-is-More (LiM) learning strategy. In addition we design a differentiable layout refinement module which maps bounding boxes to pixel-level soft masks so as to further alleviate ambiguous layout areas. Our whole model including layout prediction mask refinement and image generation can be jointly optimized in an end-to-end manner. Experimental results show that our model can generate high-quality HME images and outperforms previous generative methods. Besides a series of ablations study demonstrate effectiveness of the proposed techniques. Finally we validate that our generated images promisingly boosts the performance of HME recognition models through data augmentation. Our code and results are available at: https://github.com/AiArt-HDU/HMEG.

The recent progress in language-based open-vocabulary object detection can be largely attributed to finding better ways of leveraging large-scale data with free-form text annotations. Training such models with a discriminative objective function has proven successful but requires good positive and negative samples. However the free-form nature and the open vocabulary of object descriptions make the space of negatives extremely large. Prior works randomly sample negatives or use rule-based techniques to build them. In contrast we propose to leverage the vast knowledge built into modern generative models to automatically build negatives that are more relevant to the original data. Specifically we use large-language-models to generate negative text descriptions and text-to-image diffusion models to also generate corresponding negative images. Our experimental analysis confirms the relevance of the generated negative data and its use in language-based detectors improves performance on two complex benchmarks. Code is available at https://github.com/xiaofeng94/Gen-Enhanced-Negs.

Recovering ghost-free High Dynamic Range (HDR) images from multiple Low Dynamic Range (LDR) images becomes challenging when the LDR images exhibit saturation and significant motion. Recent Diffusion Models (DMs) have been introduced in HDR imaging field demonstrating promising performance particularly in achieving visually perceptible results compared to previous DNN-based methods. However DMs require extensive iterations with large models to estimate entire images resulting in inefficiency that hinders their practical application. To address this challenge we propose the Low-Frequency aware Diffusion (LF-Diff) model for ghost-free HDR imaging. The key idea of LF-Diff is implementing the DMs in a highly compacted latent space and integrating it into a regression-based model to enhance the details of reconstructed images. Specifically as low-frequency information is closely related to human visual perception we propose to utilize DMs to create compact low-frequency priors for the reconstruction process. In addition to take full advantage of the above low-frequency priors the Dynamic HDR Reconstruction Network (DHRNet) is carried out in a regression-based manner to obtain final HDR images. Extensive experiments conducted on synthetic and real-world benchmark datasets demonstrate that our LF-Diff performs favorably against several state-of-the-art methods and is 10x faster than previous DM-based methods.

Robotics agents often struggle to understand and follow the multi-modal prompts in complex manipulation scenes which are challenging to be sufficiently and accurately described by text alone. Moreover for long-horizon manipulation tasks the deviation from general instruction tends to accumulate if lack of intermediate guidance from high-level subgoals. For this we consider can we generate subgoal images before act to enhance the instruction following in long-horizon manipulation with multi-modal prompts? Inspired by the great success of diffusion model in image generation tasks we propose a novel hierarchical framework named as CoTDiffusion that incorporates diffusion model as a high-level planner to convert the general and multi-modal prompts into coherent visual subgoal plans which further guide the low-level policy model before action execution. We design a semantic alignment module that can anchor the progress of generated keyframes along a coherent generation chain unlocking the chain-of-thought reasoning ability of diffusion model. Additionally we propose bi-directional generation and frame concat mechanism to further enhance the fidelity of generated subgoal images and the accuracy of instruction following. The experiments cover various robotics manipulation scenarios including visual reasoning visual rearrange and visual constraints. CoTDiffusion achieves outstanding performance gain compared to the baselines without explicit subgoal generation which proves that a subgoal image is worth a thousand words of instruction.

Few-shot font generation (FFG) produces stylized font images with a limited number of reference samples which can significantly reduce labor costs in manual font designs. Most existing FFG methods follow the style-content disentanglement paradigm and employ the Generative Adversarial Network (GAN) to generate target fonts by combining the decoupled content and style representations. The complicated structure and detailed style are simultaneously generated in those methods which may be the sub-optimal solutions for FFG task. Inspired by most manual font design processes of expert designers in this paper we model font generation as a multi-stage generative process. Specifically as the injected noise and the data distribution in diffusion models can be well-separated into different sub-spaces we are able to incorporate the font transfer process into these models. Based on this observation we generalize diffusion methods to model font generative process by separating the reverse diffusion process into three stages with different functions: The structure construction stage first generates the structure information for the target character based on the source image and the font transfer stage subsequently transforms the source font to the target font. Finally the font refinement stage enhances the appearances and local details of the target font images. Based on the above multi-stage generative process we construct our font generation framework named MSD-Font with a dual-network approach to generate font images. The superior performance demonstrates the effectiveness of our model. The code is available at: https://github.com/fubinfb/MSD-Font .

We focus on the generalization ability of the 6-DoF grasp detection method in this paper. While learning-based grasp detection methods can predict grasp poses for unseen objects using the grasp distribution learned from the training set they often exhibit a significant performance drop when encountering objects with diverse shapes and structures. To enhance the grasp detection methods' generalization ability we incorporate domain prior knowledge of robotic grasping enabling better adaptation to objects with significant shape and structure differences. More specifically we employ the physical constraint regularization during the training phase to guide the model towards predicting grasps that comply with the physical rule on grasping. For the unstable grasp poses predicted on novel objects we design a contact-score joint optimization using the projection contact map to refine these poses in cluttered scenarios. Extensive experiments conducted on the GraspNet-1billion benchmark demonstrate a substantial performance gain on the novel object set and the real-world grasping experiments also demonstrate the effectiveness of our generalizing 6-DoF grasp detection method.

Event cameras capture the world at high time resolution and with minimal bandwidth requirements. However event streams which only encode changes in brightness do not contain sufficient scene information to support a wide variety of downstream tasks. In this work we design generalized event cameras that inherently preserve scene intensity in a bandwidth-efficient manner. We generalize event cameras in terms of when an event is generated and what information is transmitted. To implement our designs we turn to single-photon sensors that provide digital access to individual photon detections; this modality gives us the flexibility to realize a rich space of generalized event cameras. Our single-photon event cameras are capable of high-speed high-fidelity imaging at low readout rates. Consequently these event cameras can support plug-and-play downstream inference without capturing new event datasets or designing specialized event-vision models. As a practical implication our designs which involve lightweight and near-sensor-compatible computations provide a way to use single-photon sensors without exorbitant bandwidth costs.

Whole Slide Image (WSI) classification is often formulated as a Multiple Instance Learning (MIL) problem. Recently Vision-Language Models (VLMs) have demonstrated remarkable performance in WSI classification. However existing methods leverage coarse-grained pathogenetic descriptions for visual representation supervision which are insufficient to capture the complex visual appearance of pathogenetic images hindering the generalizability of models on diverse downstream tasks. Additionally processing high-resolution WSIs can be computationally expensive. In this paper we propose a novel "Fine-grained Visual-Semantic Interaction" (FiVE) framework for WSI classification. It is designed to enhance the model's generalizability by leveraging the interaction between localized visual patterns and fine-grained pathological semantics. Specifically with meticulously designed queries we start by utilizing a large language model to extract fine-grained pathological descriptions from various non-standardized raw reports. The output descriptions are then reconstructed into fine-grained labels used for training. By introducing a Task-specific Fine-grained Semantics (TFS) module we enable prompts to capture crucial visual information in WSIs which enhances representation learning and augments generalization capabilities significantly. Furthermore given that pathological visual patterns are redundantly distributed across tissue slices we sample a subset of visual instances during training. Our method demonstrates robust generalizability and strong transferability dominantly outperforming the counterparts on the TCGA Lung Cancer dataset with at least 9.19% higher accuracy in few-shot experiments. The code is available at: https://github.com/ls1rius/WSI_FiVE.