CVPR 2024 论文列表

Diffusion Transformer (DiT) has emerged as the new trend of generative diffusion models on image generation. In view of extremely slow convergence in typical DiT recent breakthroughs have been driven by mask strategy that significantly improves the training efficiency of DiT with additional intra-image contextual learning. Despite this progress mask strategy still suffers from two inherent limitations: (a) training-inference discrepancy and (b) fuzzy relations between mask reconstruction & generative diffusion process resulting in sub-optimal training of DiT. In this work we address these limitations by novelly unleashing the self-supervised discrimination knowledge to boost DiT training. Technically we frame our DiT in a teacher-student manner. The teacher-student discriminative pairs are built on the diffusion noises along the same Probability Flow Ordinary Differential Equation (PF-ODE). Instead of applying mask reconstruction loss over both DiT encoder and decoder we decouple DiT encoder and decoder to separately tackle discriminative and generative objectives. In particular by encoding discriminative pairs with student and teacher DiT encoders a new discriminative loss is designed to encourage the inter-image alignment in the self-supervised embedding space. After that student samples are fed into student DiT decoder to perform the typical generative diffusion task. Extensive experiments are conducted on ImageNet dataset and our method achieves a competitive balance between training cost and generative capacity.

Image diffusion models have been utilized in various tasks such as text-to-image generation and controllable image synthesis. Recent research has introduced tuning methods that make subtle adjustments to the original models yielding promising results in specific adaptations of foundational generative diffusion models. Rather than modifying the main backbone of the diffusion model we delve into the role of skip connection in U-Net and reveal that hierarchical features aggregating long-distance information across encoder and decoder make a significant impact on the content and quality of image generation. Based on the observation we propose an efficient generative tuning framework dubbed SCEdit which integrates and edits Skip Connection using a lightweight tuning module named SC-Tuner. Furthermore the proposed framework allows for straightforward extension to controllable image synthesis by injecting different conditions with Controllable SC-Tuner simplifying and unifying the network design for multi-condition inputs. Our SCEdit substantially reduces training parameters memory usage and computational expense due to its lightweight tuners with backward propagation only passing to the decoder blocks. Extensive experiments conducted on text-to-image generation and controllable image synthesis tasks demonstrate the superiority of our method in terms of efficiency and performance. Project page: https://scedit.github.io/.

This paper proposes to correct the rolling shutter (RS) distorted images by estimating the distortion flow from the global shutter (GS) to RS directly. Existing methods usually perform correction using the undistortion flow from the RS to GS. They initially predict the flow from consecutive RS frames subsequently rescaling it as the displacement fields from the RS frame to the underlying GS image using time-dependent scaling factors. Following this RS-aware forward warping is employed to convert the RS image into its GS counterpart. Nevertheless this strategy is prone to two shortcomings. First the undistortion flow estimation is rendered inaccurate by merely linear scaling the flow due to the complex non-linear motion nature. Second RS-aware forward warping often results in unavoidable artifacts. To address these limitations we introduce a new framework that directly estimates the distortion flow and rectifies the RS image with the backward warping operation. More specifically we first propose a global correlation-based flow attention mechanism to estimate the initial distortion flow and GS feature jointly which are then refined by the following coarse-to-fine decoder layers. Additionally a multi-distortion flow prediction strategy is integrated to mitigate the issue of inaccurate flow estimation further. Experimental results validate the effectiveness of the proposed method which outperforms state-of-the-art approaches on various benchmarks while maintaining high efficiency. The project is available at https://github.com/ljzycmd/DFRSC.

Feature matching is an important computer vision task that involves estimating correspondences between two images of a 3D scene and dense methods estimate all such correspondences. The aim is to learn a robust model i.e. a model able to match under challenging real-world changes. In this work we propose such a model leveraging frozen pretrained features from the foundation model DINOv2. Although these features are significantly more robust than local features trained from scratch they are inherently coarse. We therefore combine them with specialized ConvNet fine features creating a precisely localizable feature pyramid. To further improve robustness we propose a tailored transformer match decoder that predicts anchor probabilities which enables it to express multimodality. Finally we propose an improved loss formulation through regression-by-classification with subsequent robust regression. We conduct a comprehensive set of experiments that show that our method RoMa achieves significant gains setting a new state-of-the-art. In particular we achieve a 36% improvement on the extremely challenging WxBS benchmark. Code is provided at github.com/Parskatt/RoMa.

We propose RoHM an approach for robust 3D human motion reconstruction from monocular RGB(-D) videos in the presence of noise and occlusions. Most previous approaches either train neural networks to directly regress motion in 3D or learn data-driven motion priors and combine them with optimization at test time. RoHM is a novel diffusion-based motion model that conditioned on noisy and occluded input data reconstructs complete plausible motions in consistent global coordinates. Given the complexity of the problem -- requiring one to address different tasks (denoising and infilling) in different solution spaces (local and global motion) -- we decompose it into two sub-tasks and learn two models one for global trajectory and one for local motion. To capture the correlations between the two we then introduce a novel conditioning module combining it with an iterative inference scheme. We apply RoHM to a variety of tasks -- from motion reconstruction and denoising to spatial and temporal infilling. Extensive experiments on three popular datasets show that our method outperforms state-of-the-art approaches qualitatively and quantitatively while being faster at test time. The code is available at https://sanweiliti.github.io/ROHM/ROHM.html.

Lifting 2D diffusion for 3D generation is a challenging problem due to the lack of geometric prior and the complex entanglement of materials and lighting in natural images. Existing methods have shown promise by first creating the geometry through score-distillation sampling (SDS) applied to rendered surface normals followed by appearance modeling. However relying on a 2D RGB diffusion model to optimize surface normals is suboptimal due to the distribution discrepancy between natural images and normals maps leading to instability in optimization. In this paper recognizing that the normal and depth information effectively describe scene geometry and be automatically estimated from images we propose to learn a generalizable Normal-Depth diffusion model for 3D generation. We achieve this by training on the large-scale LAION dataset together with the generalizable image-to-depth and normal prior models. In an attempt to alleviate the mixed illumination effects in the generated materials we introduce an albedo diffusion model to impose data-driven constraints on the albedo component. Our experiments show that when integrated into existing text-to-3D pipelines our models significantly enhance the detail richness achieving state-of-the-art results. Our project page is at https://aigc3d.github.io/richdreamer/.

The field of image synthesis is currently flourishing due to the advancements in diffusion models. While diffusion models have been successful their computational intensity has prompted the pursuit of more efficient alternatives. As a representative work non-autoregressive Transformers (NATs) have been recognized for their rapid generation. However a major drawback of these models is their inferior performance compared to diffusion models. In this paper we aim to re-evaluate the full potential of NATs by revisiting the design of their training and inference strategies. Specifically we identify the complexities in properly configuring these strategies and indicate the possible sub-optimality in existing heuristic-driven designs. Recognizing this we propose to go beyond existing methods by directly solving the optimal strategies in an automatic framework. The resulting method named AutoNAT advances the performance boundaries of NATs notably and is able to perform comparably with the latest diffusion models with a significantly reduced inference cost. The effectiveness of AutoNAT is comprehensively validated on four benchmark datasets i.e. ImageNet-256 & 512 MS-COCO and CC3M. Code and pre-trained models will be available at https://github.com/LeapLabTHU/ImprovedNAT.

Classifier-Free Guidance (CFG) has been widely used in text-to-image diffusion models where the CFG scale is introduced to control the strength of text guidance on the whole image space. However we argue that a global CFG scale results in spatial inconsistency on varying semantic strengths and suboptimal image quality. To address this problem we present a novel approach Semantic-aware Classifier-Free Guidance (S-CFG) to customize the guidance degrees for different semantic units in text-to-image diffusion models. Specifically we first design a training-free semantic segmentation method to partition the latent image into relatively independent semantic regions at each denoising step. In particular the cross-attention map in the denoising U-net backbone is renormalized for assigning each patch to the corresponding token while the self-attention map is used to complete the semantic regions. Then to balance the amplification of diverse semantic units we adaptively adjust the CFG scales across different semantic regions to rescale the text guidance degrees into a uniform level. Finally extensive experiments demonstrate the superiority of S-CFG over the original CFG strategy on various text-to-image diffusion models without requiring any extra training cost. our codes are available at https://github.com/SmilesDZgk/S-CFG.

Recently diffusion models (DM) have been applied in magnetic resonance imaging (MRI) super-resolution (SR) reconstruction exhibiting impressive performance especially with regard to detailed reconstruction. However the current DM-based SR reconstruction methods still face the following issues: (1) They require a large number of iterations to reconstruct the final image which is inefficient and consumes a significant amount of computational resources. (2) The results reconstructed by these methods are often misaligned with the real high-resolution images leading to remarkable distortion in the reconstructed MR images. To address the aforementioned issues we propose an efficient diffusion model for multi-contrast MRI SR named as DiffMSR. Specifically we apply DM in a highly compact low-dimensional latent space to generate prior knowledge with high-frequency detail information. The highly compact latent space ensures that DM requires only a few simple iterations to produce accurate prior knowledge. In addition we design the Prior-Guide Large Window Transformer (PLWformer) as the decoder for DM which can extend the receptive field while fully utilizing the prior knowledge generated by DM to ensure that the reconstructed MR image remains undistorted. Extensive experiments on public and clinical datasets demonstrate that our DiffMSR outperforms state-of-the-art methods.

The inherent generative power of denoising diffusion models makes them well-suited for image restoration tasks where the objective is to find the optimal high-quality image within the generative space that closely resembles the input image. We propose a method to adapt a pretrained diffusion model for image restoration by simply adding noise to the input image to be restored and then denoise. Our method is based on the observation that the space of a generative model needs to be constrained. We impose this constraint by finetuning the generative model with a set of anchor images that capture the characteristics of the input image. With the constrained space we can then leverage the sampling strategy used for generation to do image restoration. We evaluate against previous methods and show superior performances on multiple real-world restoration datasets in preserving identity and image quality. We also demonstrate an important and practical application on personalized restoration where we use a personal album as the anchor images to constrain the generative space. This approach allows us to produce results that accurately preserve high-frequency details which previous works are unable to do. Project webpage: https://gen2res.github.io.

We propose residual denoising diffusion models (RDDM) a novel dual diffusion process that decouples the traditional single denoising diffusion process into residual diffusion and noise diffusion. This dual diffusion framework expands the denoising-based diffusion models initially uninterpretable for image restoration into a unified and interpretable model for both image generation and restoration by introducing residuals. Specifically our residual diffusion represents directional diffusion from the target image to the degraded input image and explicitly guides the reverse generation process for image restoration while noise diffusion represents random perturbations in the diffusion process. The residual prioritizes certainty while the noise emphasizes diversity enabling RDDM to effectively unify tasks with varying certainty or diversity requirements such as image generation and restoration. We demonstrate that our sampling process is consistent with that of DDPM and DDIM through coefficient transformation and propose a partially path-independent generation process to better understand the reverse process. Notably our RDDM enables a generic UNet trained with only an L1 loss and a batch size of 1 to compete with state-of-the-art image restoration methods. We provide code and pre-trained models to encourage further exploration application and development of our innovative framework (https://github.com/nachifur/RDDM).

Monocular depth estimation is a fundamental computer vision task. Recovering 3D depth from a single image is geometrically ill-posed and requires scene understanding so it is not surprising that the rise of deep learning has led to a breakthrough. The impressive progress of monocular depth estimators has mirrored the growth in model capacity from relatively modest CNNs to large Transformer architectures. Still monocular depth estimators tend to struggle when presented with images with unfamiliar content and layout since their knowledge of the visual world is restricted by the data seen during training and challenged by zero-shot generalization to new domains. This motivates us to explore whether the extensive priors captured in recent generative diffusion models can enable better more generalizable depth estimation. We introduce Marigold a method for affine-invariant monocular depth estimation that is derived from Stable Diffusion and retains its rich prior knowledge. The estimator can be fine-tuned in a couple of days on a single GPU using only synthetic training data. It delivers state-of-the-art performance across a wide range of datasets including over 20% performance gains in specific cases. Project page: https://marigoldmonodepth.github.io.

Portrait harmonization aims to composite a subject into a new background adjusting its lighting and color to ensure harmony with the background scene. Existing harmonization techniques often only focus on adjusting the global color and brightness of the foreground and ignore crucial illumination cues from the background such as apparent lighting direction leading to unrealistic compositions. We introduce Relightful Harmonization a lighting-aware diffusion model designed to seamlessly harmonize sophisticated lighting effect for the foreground portrait using any background image. Our approach unfolds in three stages. First we introduce a lighting representation module that allows our diffusion model to encode lighting information from target image background. Second we introduce an alignment network that aligns lighting features learned from image background with lighting features learned from panorama environment maps which is a complete representation for scene illumination. Last to further boost the photorealism of the proposed method we introduce a novel data simulation pipeline that generates synthetic training pairs from a diverse range of natural images which are used to refine the model. Our method outperforms existing benchmarks in visual fidelity and lighting coherence showing superior generalization in real-world testing scenarios highlighting its versatility and practicality.

Despite their exceptional generative abilities large T2I diffusion models much like skilled but careless artists often struggle with accurately depicting visual relationships between objects. This issue as we uncover through careful analysis arises from a misaligned text encoder that struggles to interpret specific relationships and differentiate the logical order of associated objects. To resolve this we introduce a novel task termed Relation Rectification aiming to refine the model to accurately represent a given relationship it initially fails to generate. To address this we propose an innovative solution utilizing a Heterogeneous Graph Convolutional Network (HGCN). It models the directional relationships between relation terms and corresponding objects within the input prompts. Specifically we optimize the HGCN on a pair of prompts with identical relational words but reversed object orders supplemented by a few reference images. The lightweight HGCN adjusts the text embeddings generated by the text encoder ensuring accurate reflection of the textual relation in the embedding space. Crucially our method retains the parameters of the text encoder and diffusion model preserving the model's robust performance on unrelated descriptions. We validated our approach on a newly curated dataset of diverse relational data demonstrating both quantitative and qualitative enhancements in generating images with precise visual relations. Project page: https://wuyinwei-hah.github.io/rrnet.github.io/ .

Significant advancements have been made in image editing with the recent advance of the Diffusion model. However most of the current methods primarily focus on global or subject-level modifications and often face limitations when it comes to editing specific objects when there are other objects coexisting in the scene given solely textual prompts. In response to this challenge we introduce an object-level generative task called Referring Image Editing (RIE) which enables the identification and editing of specific source objects in an image using text prompts. To tackle this task effectively we propose a tailored framework called ReferDiffusion. It aims to disentangle input prompts into multiple embeddings and employs a mixed-supervised multi-stage training strategy. To facilitate further research in this domain we introduce the RefCOCO-Edit dataset comprising images editing prompts source object segmentation masks and reference edited images for training and evaluation. Our extensive experiments demonstrate the effectiveness of our approach in identifying and editing target objects while conventional general image editing and region-based image editing methods have difficulties in this challenging task.

Image stitching from different captures often results in non-rectangular boundaries which is often considered unappealing. To solve non-rectangular boundaries current solutions involve cropping which discards image content inpainting which can introduce unrelated content or warping which can distort non-linear features and introduce artifacts. To overcome these issues we introduce a novel diffusion-based learning framework RecDiffusion for image stitching rectangling. This framework combines Motion Diffusion Models (MDM) to generate motion fields effectively transitioning from the stitched image's irregular borders to a geometrically corrected intermediary. Followed by Content Diffusion Models (CDM) for image detail refinement. Notably our sampling process utilizes a weighted map to identify regions needing correction during each iteration of CDM. Our RecDiffusion ensures geometric accuracy and overall visual appeal surpassing all previous methods in both quantitative and qualitative measures when evaluated on public benchmarks. Code is released at https://github.com/lhaippp/RecDiffusion.

We present Readout Guidance a method for controlling text-to-image diffusion models with learned signals. Readout Guidance uses readout heads lightweight networks trained to extract signals from the features of a pre-trained frozen diffusion model at every timestep. These readouts can encode single-image properties such as pose depth and edges; or higher-order properties that relate multiple images such as correspondence and appearance similarity. Furthermore by comparing the readout estimates to a user-defined target and back-propagating the gradient through the readout head these estimates can be used to guide the sampling process. Compared to prior methods for conditional generation Readout Guidance requires significantly fewer added parameters and training samples and offers a convenient and simple recipe for reproducing different forms of conditional control under a single framework with a single architecture and sampling procedure. We showcase these benefits in the applications of drag-based manipulation identity-consistent generation and spatially aligned control.

With the emergence of AR/VR 3D models are in tremendous demand. However conventional 3D modeling with Computer-Aided Design software requires much expertise and is difficult for novice users. We find that AR/VR devices in addition to serving as effective display mediums can offer a promising potential as an intuitive 3D model creation tool especially with the assistance of AI generative models. Here we propose Deep3DVRSketch the first 3D model generation network that inputs 3D VR sketches from novice users and generates highly consistent 3D models in multiple categories within seconds irrespective of the users' drawing abilities. We also contribute KO3D+ the largest 3D sketch-shape dataset. Our method pre-trains a conditional diffusion model on quality 3D data then fine-tunes an encoder to map 3D sketches onto the generator's manifold using an adaptive curriculum strategy for limited ground truths. In our experiment our approach achieves state-of-the-art performance in both model quality and fidelity with real-world input from novice users and users can even draw and obtain very detailed geometric structures. In our user study users were able to complete the 3D modeling tasks over 10 times faster using our approach compared to conventional CAD software tools. We believe that our Deep3DVRSketch and KO3D+ dataset can offer a promising solution for future 3D modeling in metaverse era. Check the project page at http://research.kokoni3d.com/Deep3DVRSketch.