CVPR 2024 论文列表

Amplitude modulated continuous-wave time-of-flight (AMCW-ToF) cameras are finding applications as flash Lidars in autonomous navigation robotics and AR/VR applications. A conventional CW-ToF camera requires illuminating the scene with a temporally varying light source and demodulating a set of quadrature measurements to recover the scene's depth and intensity. Capturing the four measurements in sequence renders the system slow invariably causing inaccuracies in depth estimates due to motion in the scene or the camera. To mitigate this problem we propose a snapshot Lidar that captures amplitude and phase simultaneously as a single time-of-flight hologram. Uniquely our approach requires minimal changes to existing CW-ToF imaging hardware. To demonstrate the efficacy of the proposed system we design and build a lab prototype and evaluate it under varying scene geometries illumination conditions and compare the reconstructed depth measurements against conventional techniques. We rigorously evaluate the robustness of our system on diverse real-world scenes to show that our technique results in a significant reduction in data bandwidth with minimal loss in reconstruction accuracy. As high-resolution CW-ToF cameras are becoming ubiquitous increasing their temporal resolution by four times enables robust real-time capture of geometries of dynamic scenes.

Contemporary models for generating images show remarkable quality and versatility. Swayed by these advantages the research community repurposes them to generate videos. Since video content is highly redundant we argue that naively bringing advances of image models to the video generation domain reduces motion fidelity visual quality and impairs scalability. In this work we build Snap Video a video-first model that systematically addresses these challenges. To do that we first extend the EDM framework to take into account spatially and temporally redundant pixels and naturally support video generation. Second we show that a U-Net--a workhorse behind image generation--scales poorly when generating videos requiring significant computational overhead. Hence we propose a new transformer-based architecture that trains 3.31 times faster than U-Nets (and is 4.5 faster at inference). This allows us to efficiently train a text-to-video model with billions of parameters for the first time reach state-of-the-art results on a number of benchmarks and generate videos with substantially higher quality temporal consistency and motion complexity. The user studies showed that our model was favored by a large margin over the most recent methods.

Recently diffusion models have made remarkable progress in text-to-image (T2I) generation synthesizing images with high fidelity and diverse contents. Despite this advancement latent space smoothness within diffusion models remains largely unexplored. Smooth latent spaces ensure that a perturbation on an input latent corresponds to a steady change in the output image. This property proves beneficial in downstream tasks including image interpolation inversion and editing. In this work we expose the non-smoothness of diffusion latent spaces by observing noticeable visual fluctuations resulting from minor latent variations. To tackle this issue we propose Smooth Diffusion a new category of diffusion models that can be simultaneously high-performing and smooth. Specifically we introduce Step-wise Variation Regularization to enforce the proportion between the variations of an arbitrary input latent and that of the output image is a constant at any diffusion training step. In addition we devise an interpolation standard deviation (ISTD) metric to effectively assess the latent space smoothness of a diffusion model. Extensive quantitative and qualitative experiments demonstrate that Smooth Diffusion stands out as a more desirable solution not only in T2I generation but also across various downstream tasks. Smooth Diffusion is implemented as a plug-and-play Smooth-LoRA to work with various community models. Code is available at https://github.com/SHI-Labs/Smooth-Diffusion.

The field of generative image inpainting and object insertion has made significant progress with the recent advent of latent diffusion models. Utilizing a precise object mask can greatly enhance these applications. However due to the challenges users encounter in creating high-fidelity masks there is a tendency for these methods to rely on more coarse masks (e.g. bounding box) for these applications. This results in limited control and compromised background content preservation. To overcome these limitations we introduce SmartMask which allows any novice user to create detailed masks for precise object insertion. Combined with a ControlNet-Inpaint model our experiments demonstrate that SmartMask achieves superior object insertion quality preserving the background content more effectively than previous methods. Notably unlike prior works the proposed approach can also be used even without user-mask guidance which allows it to perform mask-free object insertion at diverse positions and scales. Furthermore we find that when used iteratively with a novel instruction-tuning based planning model SmartMask can be used to design detailed layouts from scratch. As compared with user-scribble based layout design we observe that SmartMask allows for better quality outputs with layout-to-image generation methods.

Refractive Index Tomography is the inverse problem of reconstructing the continuously-varying 3D refractive index in a scene using 2D projected image measurements. Although a purely refractive field is not directly visible it bends light rays as they travel through space thus providing a signal for reconstruction. The effects of such fields appear in many scientific computer vision settings ranging from refraction due to transparent cells in microscopy to the lensing of distant galaxies caused by dark matter in astrophysics. Reconstructing these fields is particularly difficult due to the complex nonlinear effects of the refractive field on observed images. Furthermore while standard 3D reconstruction and tomography settings typically have access to observations of the scene from many viewpoints many refractive index tomography problem settings only have access to images observed from a single viewpoint. We introduce a method that leverages prior knowledge of light sources scattered throughout the refractive medium to help disambiguate the single-view refractive index tomography problem. We differentiably trace curved rays through a neural field representation of the refractive field and optimize its parameters to best reproduce the observed image. We demonstrate the efficacy of our approach by reconstructing simulated refractive fields analyze the effects of light source distribution on the recovered field and test our method on a simulated dark matter mapping problem where we successfully recover the 3D refractive field caused by a realistic dark matter distribution.

We introduce a framework for intrinsic latent diffusion models operating directly on the surfaces of 3D shapes with the goal of synthesizing high-quality textures. Our approach is underpinned by two contributions: Field Latents a latent representation encoding textures as discrete vector fields on the mesh vertices and Field Latent Diffusion Models which learn to denoise a diffusion process in the learned latent space on the surface. We consider a single-textured-mesh paradigm where our models are trained to generate variations of a given texture on a mesh. We show the synthesized textures are of superior fidelity compared those from existing single-textured-mesh generative models. Our models can also be adapted for user-controlled editing tasks such as inpainting and label-guided generation. The efficacy of our approach is due in part to the equivariance of our proposed framework under isometries allowing our models to seamlessly reproduce details across locally similar regions and opening the door to a notion of generative texture transfer. Code and visualizations are available at https://single-mesh-diffusion.github.io/.

While super-resolution (SR) methods based on diffusion models exhibit promising results their practical application is hindered by the substantial number of required inference steps. Recent methods utilize the degraded images in the initial state thereby shortening the Markov chain. Nevertheless these solutions either rely on a precise formulation of the degradation process or still necessitate a relatively lengthy generation path (e.g. 15 iterations). To enhance inference speed we propose a simple yet effective method for achieving single-step SR generation named SinSR. Specifically we first derive a deterministic sampling process from the most recent state-of-the-art (SOTA) method for accelerating diffusion-based SR. This allows the mapping between the input random noise and the generated high-resolution image to be obtained in a reduced and acceptable number of inference steps during training. We show that this deterministic mapping can be distilled into a student model that performs SR within only one inference step. Additionally we propose a novel consistency-preserving loss to simultaneously leverage the ground-truth image during the distillation process ensuring that the performance of the student model is not solely bound by the feature manifold of the teacher model resulting in further performance improvement. Extensive experiments conducted on synthetic and real-world datasets demonstrate that the proposed method can achieve comparable or even superior performance compared to both previous SOTA methods and the teacher model in just one sampling step resulting in a remarkable up to x10 speedup for inference. Our code will be released at https://github.com/wyf0912/SinSR/.

The recent wave of AI-generated content has witnessed the great development and success of Text-to-Image (T2I) technologies. By contrast Text-to-Video (T2V) still falls short of expectations though attracting increasing interest. Existing works either train from scratch or adapt large T2I model to videos both of which are computation and resource expensive. In this work we propose a Simple Diffusion Adapter (SimDA) that fine-tunes only 24M out of 1.1B parameters of a strong T2I model adapting it to video generation in a parameter-efficient way. In particular we turn the T2I model for T2V by designing light-weight spatial and temporal adapters for transfer learning. Besides we change the original spatial attention to the proposed Latent-Shift Attention (LSA) for temporal consistency. With a similar model architecture we further train a video super-resolution model to generate high-definition (1024 x 1024) videos. In addition to T2V generation in the wild SimDA could also be utilized in one-shot video editing with only 2 minutes tuning. Doing so our method could minimize the training effort with extremely few tunable parameters for model adaptation.

Despite the success of diffusion-based customization methods on visual content creation increasing concerns have been raised about such techniques from both privacy and political perspectives. To tackle this issue several anti-customization methods have been proposed in very recent months predominantly grounded in adversarial attacks. Unfortunately most of these methods adopt straightforward designs such as end-to-end optimization with a focus on adversarially maximizing the original training loss thereby neglecting nuanced internal properties intrinsic to the diffusion model and even leading to ineffective optimization in some diffusion time steps. In this paper we strive to bridge this gap by undertaking a comprehensive exploration of these inherent properties to boost the performance of current anti-customization approaches. Two aspects of properties are investigated: 1) We examine the relationship between time step selection and the model's perception in the frequency domain of images and find that lower time steps can give much more contributions to adversarial noises. This inspires us to propose an adaptive greedy search for optimal time steps that seamlessly integrates with existing anti-customization methods. 2) We scrutinize the roles of features at different layers during denoising and devise a sophisticated feature-based optimization framework for anti-customization. Experiments on facial benchmarks demonstrate that our approach significantly increases identity disruption thereby protecting user privacy and copyright. Our code is available at: https://github.com/somuchtome/SimAC.

Vision-Language Models (VLMs) are pretrained on large diverse and noisy web-crawled datasets. This underscores the critical need for dataset pruning as the quality of these datasets is strongly correlated with the performance of VLMs on downstream tasks. Using CLIPScore from a pretrained model to only train models using highly-aligned samples is one of the most successful methods for pruning. We argue that this approach suffers from multiple limitations including: false positives and negatives due to CLIP's pretraining on noisy labels. We propose a pruning signal Sieve that employs synthetic captions generated by image-captioning models pretrained on small diverse and well-aligned image-text pairs to evaluate the alignment of noisy image-text pairs. To bridge the gap between the limited diversity of generated captions and the high diversity of alternative text (alt-text) we estimate the semantic textual similarity in the embedding space of a language model pretrained on unlabeled text corpus. Using DataComp a multimodal dataset filtering benchmark when evaluating on 38 downstream tasks our pruning approach surpasses CLIPScore by 2.6% and 1.7% on medium and large scale respectively. In addition on retrieval tasks Sieve leads to a significant improvement of 2.7% and 4.5% on medium and large scale respectively.

A long-standing goal of 3D human reconstruction is to create lifelike and fully detailed 3D humans from single-view images. The main challenge lies in inferring unknown body shapes appearances and clothing details in areas not visible in the images. To address this we propose SiTH a novel pipeline that uniquely integrates an image-conditioned diffusion model into a 3D mesh reconstruction workflow. At the core of our method lies the decomposition of the challenging single-view reconstruction problem into generative hallucination and reconstruction subproblems. For the former we employ a powerful generative diffusion model to hallucinate unseen back-view appearance based on the input images. For the latter we leverage skinned body meshes as guidance to recover full-body texture meshes from the input and back-view images. SiTH requires as few as 500 3D human scans for training while maintaining its generality and robustness to diverse images. Extensive evaluations on two 3D human benchmarks including our newly created one highlighted our method's superior accuracy and perceptual quality in 3D textured human reconstruction. Our code and evaluation benchmark is available at https://ait.ethz.ch/sith.

Conventional image outpainting methods usually treat unobserved areas as unknown and extend the scene only in terms of semantic consistency thus overlooking the hidden information in shadows cast by unobserved areas such as the invisible shapes and semantics. In this paper we propose to extract and utilize the hidden information of unobserved areas from their shadows to enhance image outpainting. To this end we propose an end-to-end deep approach that explicitly looks into the shadows within the image. Specifically we extract shadows from the input image and identify instance-level shadow regions cast by the unobserved areas. Then the instance-level shadow representations are concatenated to predict the scene layout of each unobserved instance and outpaint the unobserved areas. Finally two discriminators are implemented to enhance alignment between the extended semantics and their shadows. In the experiments we show that our proposed approach provides complementary cues for outpainting and achieves considerable improvement on all datasets by adopting our approach as a plug-in module.

In the realm of image composition generating realistic shadow for the inserted foreground remains a formidable challenge. Previous works have developed image-to-image translation models which are trained on paired training data. However they are struggling to generate shadows with accurate shapes and intensities hindered by data scarcity and inherent task complexity. In this paper we resort to foundation model with rich prior knowledge of natural shadow images. Specifically we first adapt ControlNet to our task and then propose intensity modulation modules to improve the shadow intensity. Moreover we extend the small-scale DESOBA dataset to DESOBAv2 using a novel data acquisition pipeline. Experimental results on both DESOBA and DESOBAv2 datasets as well as real composite images demonstrate the superior capability of our model for shadow generation task. The dataset code and model are released at https://github.com/bcmi/Object-Shadow-Generation-Dataset-DESOBAv2.

Diffusion-based models have gained significant popularity for text-to-image generation due to their exceptional image-generation capabilities. A risk with these models is the potential generation of inappropriate content such as biased or harmful images. However the underlying reasons for generating such undesired content from the perspective of the diffusion model's internal representation remain unclear. Previous work interprets vectors in an interpretable latent space of diffusion models as semantic concepts. However existing approaches cannot discover directions for arbitrary concepts such as those related to inappropriate concepts. In this work we propose a novel self-supervised approach to find interpretable latent directions for a given concept. With the discovered vectors we further propose a simple approach to mitigate inappropriate generation. Extensive experiments have been conducted to verify the effectiveness of our mitigation approach namely for fair generation safe generation and responsible text-enhancing generation. Project page: https://interpretdiffusion.github.io.

In text-to-image personalization a timely and crucial challenge is the tendency of generated images overfitting to the biases present in the reference images. We initiate our study with a comprehensive categorization of the biases into background nearby-object tied-object substance (in style re-contextualization) and pose biases. These biases manifest in the generated images due to their entanglement into the subject embedding. This undesired embedding entanglement not only results in the reflection of biases from the reference images into the generated images but also notably diminishes the alignment of the generated images with the given generation prompt. To address this challenge we propose SID (Selectively Informative Description) a text description strategy that deviates from the prevalent approach of only characterizing the subject's class identification. SID is generated utilizing multimodal GPT-4 and can be seamlessly integrated into optimization-based models. We present comprehensive experimental results along with analyses of cross-attention maps subject-alignment non-subject-disentanglement and text-alignment.

Universal image restoration is a practical and potential computer vision task for real-world applications. The main challenge of this task is handling the different degradation distributions at once. Existing methods mainly utilize task-specific conditions (e.g. prompt) to guide the model to learn different distributions separately named multi-partite mapping. However it is not suitable for universal model learning as it ignores the shared information between different tasks. In this work we propose an advanced selective hourglass mapping strategy based on diffusion model termed DiffUIR. Two novel considerations make our DiffUIR non-trivial. Firstly we equip the model with strong condition guidance to obtain accurate generation direction of diffusion model (selective). More importantly DiffUIR integrates a flexible shared distribution term (SDT) into the diffusion algorithm elegantly and naturally which gradually maps different distributions into a shared one. In the reverse process combined with SDT and strong condition guidance DiffUIR iteratively guides the shared distribution to the task-specific distribution with high image quality (hourglass). Without bells and whistles by only modifying the mapping strategy we achieve state-of-the-art performance on five image restoration tasks 22 benchmarks in the universal setting and zero-shot generalization setting. Surprisingly by only using a lightweight model (only 0.89M) we could achieve outstanding performance. The source code and pre-trained models are available at https://github.com/iSEE-Laboratory/DiffUIR

We focus on a very challenging task: imaging at nighttime dynamic scenes. Most previous methods rely on the low-light enhancement of a conventional RGB camera. However they would inevitably face a dilemma between the long exposure time of nighttime and the motion blur of dynamic scenes. Event cameras react to dynamic changes with higher temporal resolution (microsecond) and higher dynamic range (120dB) offering an alternative solution. In this work we present a novel nighttime dynamic imaging method with an event camera. Specifically we discover that the event at nighttime exhibits temporal trailing characteristics and spatial non-stationary distribution. Consequently we propose a nighttime event reconstruction network (NER-Net) which mainly includes a learnable event timestamps calibration module (LETC) to align the temporal trailing events and a non-uniform illumination aware module (NIAM) to stabilize the spatiotemporal distribution of events. Moreover we construct a paired real low-light event dataset (RLED) through a co-axial imaging system including 64200 spatially and temporally aligned image GTs and low-light events. Extensive experiments demonstrate that the proposed method outperforms state-of-the-art methods in terms of visual quality and generalization ability on real-world nighttime datasets. The project are available at: https://github.com/Liu-haoyue/NER-Net.

Owe to the powerful generative priors the pre-trained text-to-image (T2I) diffusion models have become increasingly popular in solving the real-world image super-resolution problem. However as a consequence of the heavy quality degradation of input low-resolution (LR) images the destruction of local structures can lead to ambiguous image semantics. As a result the content of reproduced high-resolution image may have semantic errors deteriorating the super-resolution performance. To address this issue we present a semantics-aware approach to better preserve the semantic fidelity of generative real-world image super-resolution. First we train a degradation-aware prompt extractor which can generate accurate soft and hard semantic prompts even under strong degradation. The hard semantic prompts refer to the image tags aiming to enhance the local perception ability of the T2I model while the soft semantic prompts compensate for the hard ones to provide additional representation information. These semantic prompts encourage the T2I model to generate detailed and semantically accurate results. Furthermore during the inference process we integrate the LR images into the initial sampling noise to mitigate the diffusion model's tendency to generate excessive random details. The experiments show that our method can reproduce more realistic image details and hold better the semantics. The source code of our method can be found at https://github.com/cswry/SeeSR