CVPR 2024 论文列表

Diffusion models have made tremendous progress in text-driven image and video generation. Now text-to-image foundation models are widely applied to various downstream image synthesis tasks such as controllable image generation and image editing while downstream video synthesis tasks are less explored for several reasons. First it requires huge memory and computation overhead to train a video generation foundation model. Even with video foundation models additional costly training is still required for downstream video synthesis tasks. Second although some works extend image diffusion models into videos in a training-free manner temporal consistency cannot be well preserved. Finally these adaption methods are specifically designed for one task and fail to generalize to different tasks. To mitigate these issues we propose a training-free general-purpose video synthesis framework coined as BIVDiff via bridging specific image diffusion models and general text-to-video foundation diffusion models. Specifically we first use a specific image diffusion model (e.g. ControlNet and Instruct Pix2Pix) for frame-wise video generation then perform Mixed Inversion on the generated video and finally input the inverted latents into the video diffusion models (e.g. VidRD and ZeroScope) for temporal smoothing. This decoupled framework enables flexible image model selection for different purposes with strong task generalization and high efficiency. To validate the effectiveness and general use of BIVDiff we perform a wide range of video synthesis tasks including controllable video generation video editing video inpainting and outpainting.

Diffusion models (DMs) have exhibited superior performance in generating high-quality and diverse images. However this exceptional performance comes at the cost of expensive generation process particularly due to the heavily used attention module in leading models. Existing works mainly adopt a retraining process to enhance DM efficiency. This is computationally expensive and not very scalable. To this end we introduce the Attention-driven Training-free Efficient Diffusion Model (AT-EDM) framework that leverages attention maps to perform run-time pruning of redundant tokens without the need for any retraining. Specifically for single-denoising-step pruning we develop a novel ranking algorithm Generalized Weighted Page Rank (G-WPR) to identify redundant tokens and a similarity-based recovery method to restore tokens for the convolution operation. In addition we propose a Denoising-Steps-Aware Pruning (DSAP) approach to adjust the pruning budget across different denoising timesteps for better generation quality. Extensive evaluations show that AT-EDM performs favorably against prior art in terms of efficiency (e.g. 38.8% FLOPs saving and up to 1.53x speed-up over Stable Diffusion XL) while maintaining nearly the same FID and CLIP scores as the full model.

Monocular depth estimation has experienced significant progress on terrestrial images in recent years thanks to deep learning advancements. But it remains inadequate for underwater scenes primarily due to data scarcity. Given the inherent challenges of light attenuation and backscatter in water acquiring clear underwater images or precise depth is notably difficult and costly. To mitigate this issue learning-based approaches often rely on synthetic data or turn to self- or unsupervised manners. Nonetheless their performance is often hindered by domain gap and looser constraints. In this paper we propose a novel pipeline for generating photorealistic underwater images using accurate terrestrial depth. This approach facilitates the supervised training of models for underwater depth estimation effectively reducing the performance disparity between terrestrial and underwater environments. Contrary to previous synthetic datasets that merely apply style transfer to terrestrial images without scene content change our approach uniquely creates vivid non-existent underwater scenes by leveraging terrestrial depth data through the innovative Stable Diffusion model. Specifically we introduce a specialized Depth2Underwater ControlNet trained on prepared \ Underwater Depth Text\ data triplets for this generation task. Our newly developed dataset Atlantis enables terrestrial depth estimation models to achieve considerable improvements on unseen underwater scenes surpassing their terrestrial pretrained counterparts both quantitatively and qualitatively. Moreover we further show its practical utility by applying the improved depth in underwater image enhancement and its smaller domain gap from the LLVM perspective. Code and dataset are publicly available at https://github.com/zkawfanx/Atlantis.

We present As-Plausible-as-Possible (APAP) mesh deformation technique that leverages 2D diffusion priors to preserve the plausibility of a mesh under user-controlled deformation. Our framework uses per-face Jacobians to represent mesh deformations where mesh vertex coordinates are computed via a differentiable Poisson Solve. The deformed mesh is rendered and the resulting 2D image is used in the Score Distillation Sampling (SDS) process which enables extracting meaningful plausibility priors from a pretrained 2D diffusion model. To better preserve the identity of the edited mesh we fine-tune our 2D diffusion model with LoRA. Gradients extracted by SDS and a user-prescribed handle displacement are then backpropagated to the per-face Jacobians and we use iterative gradient descent to compute the final deformation that balances between the user edit and the output plausibility. We evaluate our method with 2D and 3D meshes and demonstrate qualitative and quantitative improvements when using plausibility priors over geometry-preservation or distortion-minimization priors used by previous techniques. Our project page is at: https://as-plausible-aspossible.github.io/

Super-resolution (SR) and image generation are important tasks in computer vision and are widely adopted in real-world applications. Most existing methods however generate images only at fixed-scale magnification and suffer from over-smoothing and artifacts. Additionally they do not offer enough diversity of output images nor image consistency at different scales. Most relevant work applied Implicit Neural Representation (INR) to the denoising diffusion model to obtain continuous-resolution yet diverse and high-quality SR results. Since this model operates in the image space the larger the resolution of image is produced the more memory and inference time is required and it also does not maintain scale-specific consistency. We propose a novel pipeline that can super-resolve an input image or generate from a random noise a novel image at arbitrary scales. The method consists of a pretrained auto-encoder a latent diffusion model and an implicit neural decoder and their learning strategies. The proposed method adopts diffusion processes in a latent space thus efficient yet aligned with output image space decoded by MLPs at arbitrary scales. More specifically our arbitrary-scale decoder is designed by the symmetric decoder w/o up-scaling from the pretrained auto-encoder and Local Implicit Image Function (LIIF) in series. The latent diffusion process is learnt by the denoising and the alignment losses jointly. Errors in output images are backpropagated via the fixed decoder improving the quality of output images. In the extensive experiments using multiple public benchmarks on the two tasks i.e. image super-resolution and novel image generation at arbitrary scales the proposed method outperforms relevant methods in metrics of image quality diversity and scale consistency. It is significantly better than the relevant prior-art in the inference speed and memory usage.

Computer animation's quest to bridge content and style has historically been a challenging venture with previous efforts often leaning toward one at the expense of the other. This paper tackles the inherent challenge of content-style duality ensuring a harmonious fusion where the core narrative of the content is both preserved and elevated through stylistic enhancements. We propose a novel Multi-condition Motion Latent Diffusion Model (MCM-LDM) for Arbitrary Motion Style Transfer (AMST). Our MCM-LDM significantly emphasizes preserving trajectories recognizing their fundamental role in defining the essence and fluidity of motion content. Our MCM-LDM's cornerstone lies in its ability first to disentangle and then intricately weave together motion's tripartite components: motion trajectory motion content and motion style. The critical insight of MCM-LDM is to embed multiple conditions with distinct priorities. The content channel serves as the primary flow guiding the overall structure and movement while the trajectory and style channels act as auxiliary components and synchronize with the primary one dynamically. This mechanism ensures that multi-conditions can seamlessly integrate into the main flow enhancing the overall animation without overshadowing the core content. Empirical evaluations underscore the model's proficiency in achieving fluid and authentic motion style transfers setting a new benchmark in the realm of computer animation. The source code and model are available at https://github.com/XingliangJin/MCM-LDM.git.

Recent advancements in text-to-image technology have significantly advanced the field of image customization. Among various applications the task of customizing diverse scenes for user-specified composited elements holds great application value but has not been extensively explored. Addressing this gap we propose AnyScene a specialized framework designed to create varied scenes from composited foreground using textual prompts. AnyScene addresses the primary challenges inherent in existing methods particularly scene disharmony due to a lack of foreground semantic understanding and distortion of foreground elements. Specifically we develop a foreground injection module that guides a pre-trained diffusion model to generate cohesive scenes in visual harmony with the provided foreground. To enhance robust generation we implement a layout control strategy that prevents distortions of foreground elements. Furthermore an efficient image blending mechanism seamlessly reintegrates foreground details into the generated scenes producing outputs with overall visual harmony and precise foreground details. In addition we propose a new benchmark and a series of quantitative metrics to evaluate this proposed image customization task. Extensive experimental results demonstrate the effectiveness of AnyScene which confirms its potential in various applications.

With the advancement of generative models the assessment of generated images becomes increasingly more important. Previous methods measure distances between features of reference and generated images from trained vision models. In this paper we conduct an extensive investigation into the relationship between the representation space and input space around generated images. We first propose two measures related to the presence of unnatural elements within images: complexity which indicates how non-linear the representation space is and vulnerability which is related to how easily the extracted feature changes by adversarial input changes. Based on these we introduce a new metric to evaluating image-generative models called anomaly score (AS). Moreover we propose AS-i (anomaly score for individual images) that can effectively evaluate generated images individually. Experimental results demonstrate the validity of the proposed approach.

Character Animation aims to generating character videos from still images through driving signals. Currently diffusion models have become the mainstream in visual generation research owing to their robust generative capabilities. However challenges persist in the realm of image-to-video especially in character animation where temporally maintaining consistency with detailed information from character remains a formidable problem. In this paper we leverage the power of diffusion models and propose a novel framework tailored for character animation. To preserve consistency of intricate appearance features from reference image we design ReferenceNet to merge detail features via spatial attention. To ensure controllability and continuity we introduce an efficient pose guider to direct character's movements and employ an effective temporal modeling approach to ensure smooth inter-frame transitions between video frames. By expanding the training data our approach can animate arbitrary characters yielding superior results in character animation compared to other image-to-video methods. Furthermore we evaluate our method on image animation benchmarks achieving state-of-the-art results.

Diffusion models currently dominate the field of data-driven image synthesis with their unparalleled scaling to large datasets. In this paper we identify and rectify several causes for uneven and ineffective training in the popular ADM diffusion model architecture without altering its high-level structure. Observing uncontrolled magnitude changes and imbalances in both the network activations and weights over the course of training we redesign the network layers to preserve activation weight and update magnitudes on expectation. We find that systematic application of this philosophy eliminates the observed drifts and imbalances resulting in considerably better networks at equal computational complexity. Our modifications improve the previous record FID of 2.41 in ImageNet-512 synthesis to 1.81 achieved using fast deterministic sampling. As an independent contribution we present a method for setting the exponential moving average (EMA) parameters post-hoc i.e. after completing the training run. This allows precise tuning of EMA length without the cost of performing several training runs and reveals its surprising interactions with network architecture training time and guidance.

Denoising diffusion probabilistic models (DDPMs) employ a sequence of white Gaussian noise samples to generate an image. In analogy with GANs those noise maps could be considered as the latent code associated with the generated image. However this native noise space does not possess a convenient structure and is thus challenging to work with in editing tasks. Here we propose an alternative latent noise space for DDPM that enables a wide range of editing operations via simple means and present an inversion method for extracting these edit-friendly noise maps for any given image (real or synthetically generated). As opposed to the native DDPM noise space the edit-friendly noise maps do not have a standard normal distribution and are not statistically independent across timesteps. However they allow perfect reconstruction of any desired image and simple transformations on them translate into meaningful manipulations of the output image (e.g. shifting color edits). Moreover in text-conditional models fixing those noise maps while changing the text prompt modifies semantics while retaining structure. We illustrate how this property enables text-based editing of real images via the diverse DDPM sampling scheme (in contrast to the popular non-diverse DDIM inversion). We also show how it can be used within existing diffusion-based editing methods to improve their quality and diversity. The code of the method is attached to this submission.

Text-guided diffusion models have revolutionized image and video generation and have also been successfully used for optimization-based 3D object synthesis. Here we instead focus on the underexplored text-to-4D setting and synthesize dynamic animated 3D objects using score distillation methods with an additional temporal dimension. Compared to previous work we pursue a novel compositional generation-based approach and combine text-to-image text-to-video and 3D-aware multiview diffusion models to provide feedback during 4D object optimization thereby simultaneously enforcing temporal consistency high-quality visual appearance and realistic geometry. Our method called Align Your Gaussians (AYG) leverages dynamic 3D Gaussian Splatting with deformation fields as 4D representation. Crucial to AYG is a novel method to regularize the distribution of the moving 3D Gaussians and thereby stabilize the optimization and induce motion. We also propose a motion amplification mechanism as well as a new autoregressive synthesis scheme to generate and combine multiple 4D sequences for longer generation. These techniques allow us to synthesize vivid dynamic scenes outperform previous work qualitatively and quantitatively and achieve state-of-the-art text-to-4D performance. Due to the Gaussian 4D representation different 4D animations can be seamlessly combined as we demonstrate. AYG opens up promising avenues for animation simulation and digital content creation as well as synthetic data generation.

We propose a method to control material attributes of objects like roughness metallic albedo and transparency in real images. Our method capitalizes on the generative prior of text-to-image models known for photorealism employing a scalar value and instructions to alter low-level material properties. Addressing the lack of datasets with controlled material attributes we generated an object-centric synthetic dataset with physically-based materials. Fine-tuning a modified pre-trained text-to-image model on this synthetic dataset enables us to edit material properties in real-world images while preserving all other attributes. We show the potential application of our model to material edited NeRFs.

Existing score distillation methods are sensitive to classifier-free guidance (CFG) scale manifested as over-smoothness or instability at small CFG scales while over-saturation at large ones. To explain and analyze these issues we revisit the derivation of Score Distillation Sampling (SDS) and decipher existing score distillation with the Wasserstein Generative Adversarial Network (WGAN) paradigm. With the WGAN paradigm we find that existing score distillation either employs a fixed sub-optimal discriminator or conducts incomplete discriminator optimization resulting in the scale-sensitive issue. We propose the Adversarial Score Distillation (ASD) which maintains an optimizable discriminator and updates it using the complete optimization objective. Experiments show that the proposed ASD performs favorably in 2D distillation and text-to-3D tasks against existing methods. Furthermore to explore the generalization ability of our paradigm we extend ASD to the image editing task which achieves competitive results. The project page and code are at https://github.com/2y7c3/ASD

Action detection aims to localize the starting and ending points of action instances in untrimmed videos and predict the classes of those instances. In this paper we make the observation that the outputs of the action detection task can be formulated as images. Thus from a novel perspective we tackle action detection via a three-image generation process to generate starting point ending point and action-class predictions as images via our proposed Action Detection Image Diffusion (ADI-Diff) framework. Furthermore since our images differ from natural images and exhibit special properties we further explore a Discrete Action-Detection Diffusion Process and a Row-Column Transformer design to better handle their processing. Our ADI-Diff framework achieves state-of-the-art results on two widely-used datasets.

Diffusion probabilistic models (DPMs) have shown remarkable performance in high-resolution image synthesis but their sampling efficiency is still to be desired due to the typically large number of sampling steps. Recent advancements in high-order numerical ODE solvers for DPMs have enabled the generation of high-quality images with much fewer sampling steps. While this is a significant development most sampling methods still employ uniform time steps which is not optimal when using a small number of steps. To address this issue we propose a general framework for designing an optimization problem that seeks more appropriate time steps for a specific numerical ODE solver for DPMs. This optimization problem aims to minimize the distance between the ground-truth solution to the ODE and an approximate solution corresponding to the numerical solver. It can be efficiently solved using the constrained trust region method taking less than 15 seconds. Our extensive experiments on both unconditional and conditional sampling using pixel- and latent-space DPMs demonstrate that when combined with the state-of-the-art sampling method UniPC our optimized time steps significantly improve image generation performance in terms of FID scores for datasets such as CIFAR-10 and ImageNet compared to using uniform time steps.

We present MM-AU a novel dataset for Multi-Modal Accident video Understanding. MM-AU contains 11727 in-the-wild ego-view accident videos each with temporally aligned text descriptions. We annotate over 2.23 million object boxes and 58650 pairs of video-based accident reasons covering 58 accident categories. MM-AU supports various accident understanding tasks particularly multimodal video diffusion to understand accident cause-effect chains for safe driving. With MM-AU we present an Abductive accident Video understanding framework for Safe Driving perception (AdVersa-SD). AdVersa-SD performs video diffusion via an Object-Centric Video Diffusion (OAVD) method which is driven by an abductive CLIP model. This model involves a contrastive interaction loss to learn the pair co-occurrence of normal near-accident accident frames with the corresponding text descriptions such as accident reasons prevention advice and accident categories. OAVD enforces the object region learning while fixing the content of the original frame background in video generation to find the dominant objects for certain accidents. Extensive experiments verify the abductive ability of AdVersa-SD and the superiority of OAVD against the state-of-the-art diffusion models. Additionally we provide careful benchmark evaluations for object detection and accident reason answering since AdVersa-SD relies on precise object and accident reason information.

Recent advances in diffusion models have successfully enabled text-guided image inpainting. While it seems straightforward to extend such editing capability into the video domain there have been fewer works regarding text-guided video inpainting. Given a video a masked region at its initial frame and an editing prompt it requires a model to do infilling at each frame following the editing guidance while keeping the out-of-mask region intact. There are three main challenges in text-guided video inpainting: (i) temporal consistency of the edited video (ii) supporting different inpainting types at different structural fidelity levels and (iii) dealing with variable video length. To address these challenges we introduce Any-Length Video Inpainting with Diffusion Model dubbed as AVID. At its core our model is equipped with effective motion modules and adjustable structure guidance for fixed-length video inpainting. Building on top of that we propose a novel Temporal MultiDiffusion sampling pipeline with a middle-frame attention guidance mechanism facilitating the generation of videos with any desired duration. Our comprehensive experiments show our model can robustly deal with various inpainting types at different video duration ranges with high quality.