CVPR 2024 论文列表

This paper proposes ConsistDreamer - a novel framework that lifts 2D diffusion models with 3D awareness and 3D consistency thus enabling high-fidelity instruction-guided scene editing. To overcome the fundamental limitation of missing 3D consistency in 2D diffusion models our key insight is to introduce three synergetic strategies that augment the input of the 2D diffusion model to become 3D-aware and to explicitly enforce 3D consistency during the training process. Specifically we design surrounding views as context-rich input for the 2D diffusion model and generate 3D-consistent structured noise instead of image-independent noise. Moreover we introduce self-supervised consistency-enforcing training within the per-scene editing procedure. Extensive evaluation shows that our ConsistDreamer achieves state-of-the-art performance for instruction-guided scene editing across various scenes and editing instructions particularly in complicated large-scale indoor scenes from ScanNet++ with significantly improved sharpness and fine-grained textures. Notably ConsistDreamer stands as the first work capable of successfully editing complex (e.g. plaid/checkered) patterns.

Addressing pose ambiguity in 6D object pose estimation from single RGB images presents a significant challenge particularly due to object symmetries or occlusions. In response we introduce a novel score-based diffusion method applied to the SE(3) group marking the first application of diffusion models to SE(3) within the image domain specifically tailored for pose estimation tasks. Extensive evaluations demonstrate the method's efficacy in handling pose ambiguity mitigating perspective-induced ambiguity and showcasing the robustness of our surrogate Stein score formulation on SE(3). This formulation not only improves the convergence of denoising process but also enhances computational efficiency. Thus we pioneer a promising strategy for 6D object pose estimation.

While there has been significant progress in customizing text-to-image generation models generating images that combine multiple personalized concepts remains challenging. In this work we introduce Concept Weaver a method for composing customized text-to-image diffusion models at inference time. Specifically the method breaks the process into two steps: creating a template image aligned with the semantics of input prompts and then personalizing the template using a concept fusion strategy. The fusion strategy incorporates the appearance of the target concepts into the template image while retaining its structural details. The results indicate that our method can generate multiple custom concepts with higher identity fidelity compared to alternative approaches. Furthermore the method is shown to seamlessly handle more than two concepts and closely follow the semantic meaning of the input prompt without blending appearances across different subjects.

In this work we propose a method to address the challenge of rendering a 3D human from a single image in a free-view manner. Some existing approaches could achieve this by using generalizable pixel-aligned implicit fields to reconstruct a textured mesh of a human or by employing a 2D diffusion model as guidance with the Score Distillation Sampling (SDS) method to lift the 2D image into 3D space. However a generalizable implicit field often results in an over-smooth texture field while the SDS method tends to lead to a texture-inconsistent novel view with the input image. In this paper we introduce a texture-consistent back view synthesis method that could transfer the reference image content to the back view through depth-guided mutual self-attention. With this method we could achieve high-fidelity and texture-consistent human rendering from a single image. Moreover to alleviate the color distortion that occurs in the side region we propose a visibility-aware patch consistency regularization combined with the synthesized back view texture. Experiments conducted on both real and synthetic data demonstrate the effectiveness of our method and show that our approach outperforms previous baseline methods.

We train a set of open text-to-image (T2I) diffusion models on a dataset of curated Creative-Commons-licensed (CC) images which yields models that are competitive with Stable Diffusion 2 (SD2). This task presents two challenges: (1) high-resolution CC images lack the captions necessary to train T2I models; (2) CC images are relatively scarce. To address these challenges we use an intuitive transfer learning technique to produce a set of high-quality synthetic captions paired with our assembled CC images. We then develop a data- and compute-efficient training recipe that requires as little as 3% of the LAION data (i.e. roughly 70 million examples) needed to train existing SD2 models but obtains the same quality. These results indicate that we have a sufficient number of CC images (also roughly 70 million) for training high-quality models. Our recipe also implements a variety of optimizations that achieve 2.71x training speed-ups enabling rapid model iteration. We leverage this recipe to train several high-quality T2I mod- els which we dub the CommonCanvas family. Our largest model achieves comparable performance to SD2 on human evaluation even though we use a synthetically captioned CC-image dataset that is only <3% the size of LAION for training. We release our models data and code on GitHub.

Domain Generalized Semantic Segmentation (DGSS) deals with training a model on a labeled source domain with the aim of generalizing to unseen domains during inference. Existing DGSS methods typically effectuate robust features by means of Domain Randomization (DR). Such an approach is often limited as it can only account for style diversification and not content. In this work we take an orthogonal approach to DGSS and propose to use an assembly of CoLlaborative FOUndation models for Domain Generalized Semantic Segmentation (CLOUDS). In detail CLOUDS is a framework that integrates Foundation Models of various kinds: (i) CLIP backbone for its robust feature representation (ii) Diffusion Model to diversify the content thereby covering various modes of the possible target distribution and (iii) Segment Anything Model (SAM) for iteratively refining the predictions of the segmentation model. Extensive experiments show that our CLOUDS excels in adapting from synthetic to real DGSS benchmarks and under varying weather conditions notably outperforming prior methods by 5.6% and 6.7% on averaged mIoU respectively. Our code is available at https://github.com/yasserben/CLOUDS

Diffusion model is a promising approach to image generation and has been employed for Pose-Guided Person Image Synthesis (PGPIS) with competitive performance. While existing methods simply align the person appearance to the target pose they are prone to overfitting due to the lack of a high-level semantic understanding on the source person image. In this paper we propose a novel Coarse-to-Fine Latent Diffusion (CFLD) method for PGPIS. In the absence of image-caption pairs and textual prompts we develop a novel training paradigm purely based on images to control the generation process of a pre-trained text-to-image diffusion model. A perception-refined decoder is designed to progressively refine a set of learnable queries and extract semantic understanding of person images as a coarse-grained prompt. This allows for the decoupling of fine-grained appearance and pose information controls at different stages and thus circumventing the potential overfitting problem. To generate more realistic texture details a hybrid-granularity attention module is proposed to encode multi-scale fine-grained appearance features as bias terms to augment the coarse-grained prompt. Both quantitative and qualitative experimental results on the DeepFashion benchmark demonstrate the superiority of our method over the state of the arts for PGPIS. Code is available at https://github.com/YanzuoLu/CFLD.

Existing super-resolution (SR) models primarily focus on restoring local texture details often neglecting the global semantic information within the scene. This oversight can lead to the omission of crucial semantic details or the introduction of inaccurate textures during the recovery process. In our work we introduce the Cognitive Super-Resolution (CoSeR) framework empowering SR models with the capacity to comprehend low-resolution images. We achieve this by marrying image appearance and language understanding to generate a cognitive embedding which not only activates prior information from large text-to-image diffusion models but also facilitates the generation of high-quality reference images to optimize the SR process. To further improve image fidelity we propose a novel condition injection scheme called "All-in-Attention" consolidating all conditional information into a single module. Consequently our method successfully restores semantically correct and photorealistic details demonstrating state-of-the-art performance across multiple benchmarks. Project page: https://coser-main.github.io/

Large generative diffusion models have revolutionized text-to-image generation and offer immense potential for conditional generation tasks such as image enhancement restoration editing and compositing. However their widespread adoption is hindered by the high computational cost which limits their real-time application. To address this challenge we introduce a novel method dubbed CoDi that adapts a pre-trained latent diffusion model to accept additional image conditioning inputs while significantly reducing the sampling steps required to achieve high-quality results. Our method can leverage architectures such as ControlNet to incorporate conditioning inputs without compromising the model's prior knowledge gained during large scale pre-training. Additionally a conditional consistency loss enforces consistent predictions across diffusion steps effectively compelling the model to generate high-quality images with conditions in a few steps. Our conditional-task learning and distillation approach outperforms previous distillation methods achieving a new state-of-the-art in producing high-quality images with very few steps (e.g. 1-4) across multiple tasks including super-resolution text-guided image editing and depth-to-image generation.

Co-speech gestures if presented in the lively form of videos can achieve superior visual effects in human-machine interaction. While previous works mostly generate structural human skeletons resulting in the omission of appearance information we focus on the direct generation of audio-driven co-speech gesture videos in this work. There are two main challenges: 1) A suitable motion feature is needed to describe complex human movements with crucial appearance information. 2) Gestures and speech exhibit inherent dependencies and should be temporally aligned even of arbitrary length. To solve these problems we present a novel motion-decoupled framework to generate co-speech gesture videos. Specifically we first introduce a well-designed nonlinear TPS transformation to obtain latent motion features preserving essential appearance information. Then a transformer-based diffusion model is proposed to learn the temporal correlation between gestures and speech and performs generation in the latent motion space followed by an optimal motion selection module to produce long-term coherent and consistent gesture videos. For better visual perception we further design a refinement network focusing on missing details of certain areas. Extensive experimental results show that our proposed framework significantly outperforms existing approaches in both motion and video-related evaluations. Our code demos and more resources are available at https://github.com/thuhcsi/S2G-MDDiffusion.

Existing multi-person human reconstruction approaches mainly focus on recovering accurate poses or avoiding penetration but overlook the modeling of close interactions. In this work we tackle the task of reconstructing closely interactive humans from a monocular video. The main challenge of this task comes from insufficient visual information caused by depth ambiguity and severe inter-person occlusion. In view of this we propose to leverage knowledge from proxemic behavior and physics to compensate the lack of visual information. This is based on the observation that human interaction has specific patterns following the social proxemics. Specifically we first design a latent representation based on Vector Quantised-Variational AutoEncoder (VQ-VAE) to model human interaction. A proxemics and physics guided diffusion model is then introduced to denoise the initial distribution. We design the diffusion model as dual branch with each branch representing one individual such that the interaction can be modeled via cross attention. With the learned priors of VQ-VAE and physical constraint as the additional information our proposed approach is capable of estimating accurate poses that are also proxemics and physics plausible. Experimental results on Hi4D 3DPW and CHI3D demonstrate that our method outperforms existing approaches. The code is available at https://github.com/boycehbz/HumanInteraction.

This work aims to improve the efficiency of text-to-image diffusion models. While diffusion models use computationally expensive UNet-based denoising operations in every generation step we identify that not all operations are equally relevant for the final output quality. In particular we observe that UNet layers operating on high-res feature maps are relatively sensitive to small perturbations. In contrast low-res feature maps influence the semantic layout of the final image and can often be perturbed with no noticeable change in the output. Based on this observation we propose Clockwork Diffusion a method that periodically reuses computation from preceding denoising steps to approximate low-res feature maps at one or more subsequent steps. For multiple base- lines and for both text-to-image generation and image editing we demonstrate that Clockwork leads to comparable or improved perceptual scores with drastically reduced computational complexity. As an example for Stable Diffusion v1.5 with 8 DPM++ steps we save 32% of FLOPs with negligible FID and CLIP change. We re- lease code at https://github.com/Qualcomm-AI-research/clockwork-diffusion

3D city generation is a desirable yet challenging task since humans are more sensitive to structural distortions in urban environments. Additionally generating 3D cities is more complex than 3D natural scenes since buildings as objects of the same class exhibit a wider range of appearances compared to the relatively consistent appearance of objects like trees in natural scenes. To address these challenges we propose CityDreamer a compositional generative model designed specifically for unbounded 3D cities. Our key insight is that 3D city generation should be a composition of different types of neural fields: 1) various building instances and 2) background stuff such as roads and green lands. Specifically we adopt the bird's eye view scene representation and employ a volumetric render for both instance-oriented and stuff-oriented neural fields. The generative hash grid and periodic positional embedding are tailored as scene parameterization to suit the distinct characteristics of building instances and background stuff. Furthermore we contribute a suite of CityGen Datasets including OSM and GoogleEarth which comprises a vast amount of real-world city imagery to enhance the realism of the generated 3D cities both in their layouts and appearances. CityDreamer achieves state-of-the-art performance not only in generating realistic 3D cities but also in localized editing within the generated cities.

Diffusion models have recently achieved remarkable progress in generating realistic images. However challenges remain in accurately understanding and synthesizing the layout requirements in the textual prompts. To align the generated image with layout instructions we present a training-free layout calibration system SimM that intervenes in the generative process on the fly during inference time. Specifically following a "check-locate-rectify" pipeline the system first analyses the prompt to generate the target layout and compares it with the intermediate outputs to automatically detect errors. Then by moving the located activations and making intra- and inter-map adjustments the rectification process can be performed with negligible computational overhead. To evaluate SimM over a range of layout requirements we present a benchmark SimMBench that compensates for the lack of superlative spatial relations in existing datasets. And both quantitative and qualitative results demonstrate the effectiveness of the proposed SimM in calibrating the layout inconsistencies. Our project page is at https://simm-t2i.github.io/SimM.

Multi-view diffusion models obtained by applying Supervised Finetuning (SFT) to text-to-image diffusion models have driven recent breakthroughs in text-to-3D research. However due to the limited size and quality of existing 3D datasets they still suffer from multi-view inconsistencies and Neural Radiance Field (NeRF) reconstruction artifacts. We argue that multi-view diffusion models can benefit from further Reinforcement Learning Finetuning (RLFT) which allows models to learn from the data generated by themselves and improve beyond their dataset limitations during SFT. To this end we introduce Carve3D an improved RLFT algorithm coupled with a novel Multi-view Reconstruction Consistency (MRC) metric to enhance the consistency of multi-view diffusion models. To measure the MRC metric on a set of multi-view images we compare them with their corresponding NeRF renderings at the same camera viewpoints. The resulting model which we denote as Carve3DM demonstrates superior multi-view consistency and NeRF reconstruction quality than existing models. Our results suggest that pairing SFT with Carve3D's RLFT is essential for developing multi-view-consistent diffusion models mirroring the standard Large Language Model (LLM) alignment pipeline. Our code training and testing data and video results are available at: https://desaixie.github.io/carve-3d.

We concentrate on a novel human-centric image synthesis task that is given only one reference facial photograph it is expected to generate specific individual images with diverse head positions poses facial expressions and illuminations in different contexts. To accomplish this goal we argue that our generative model should be capable of the following favorable characteristics: (1) a strong visual and semantic understanding of our world and human society for basic object and human image generation. (2) generalizable identity preservation ability. (3) flexible and fine-grained head control. Recently large pre-trained text-to-image diffusion models have shown remarkable results serving as a powerful generative foundation. As a basis we aim to unleash the above two capabilities of the pre-trained model. In this work we present a new framework named CapHuman. We embrace the "encode then learn to align" paradigm which enables generalizable identity preservation for new individuals without cumbersome tuning at inference. CapHuman encodes identity features and then learns to align them into the latent space. Moreover we introduce the 3D facial prior to equip our model with control over the human head in a flexible and 3D-consistent manner. Extensive qualitative and quantitative analyses demonstrate our CapHuman can produce well-identity-preserved photo-realistic and high-fidelity portraits with content-rich representations and various head renditions superior to established baselines. Code and checkpoint will be released at https://github.com/VamosC/CapHuman.

Stable Diffusion has established itself as a foundation model in generative AI artistic applications receiving widespread research and application. Some recent fine-tuning methods have made it feasible for individuals to implant personalized concepts onto the basic Stable Diffusion model with minimal computational costs on small datasets. However these innovations have also given rise to issues like facial privacy forgery and artistic copyright infringement. In recent studies researchers have explored the addition of imperceptible adversarial perturbations to images to prevent potential unauthorized exploitation and infringements when personal data is used for fine-tuning Stable Diffusion. Although these studies have demonstrated the ability to protect images it is essential to consider that these methods may not be entirely applicable in real-world scenarios. In this paper we systematically evaluate the use of perturbations to protect images within a practical threat model. The results suggest that these approaches may not be sufficient to safeguard image privacy and copyright effectively. Furthermore we introduce a purification method capable of removing protected perturbations while preserving the original image structure to the greatest extent possible. Experiments reveal that Stable Diffusion can effectively learn from purified images over all protective methods.

Diffusion models have recently revolutionized the field of image synthesis due to their ability to generate photorealistic images. However one of the major drawbacks of diffusion models is that the image generation process is costly. A large image-to-image network has to be applied many times to iteratively refine an image from random noise. While many recent works propose techniques to reduce the number of required steps they generally treat the underlying denoising network as a black box. In this work we investigate the behavior of the layers within the network and find that 1) the layers' output changes smoothly over time 2) the layers show distinct patterns of change and 3) the change from step to step is often very small. We hypothesize that many layer computations in the denoising network are redundant. Leveraging this we introduce Block Caching in which we reuse outputs from layer blocks of previous steps to speed up inference. Furthermore we propose a technique to automatically determine caching schedules based on each block's changes over timesteps. In our experiments we show through FID human evaluation and qualitative analysis that Block Caching allows to generate images with higher visual quality at the same computational cost. We demonstrate this for different state-of-the-art models (LDM and EMU) and solvers (DDIM and DPM).

To achieve greater accuracy hypergraph matching algorithms require exponential increases in computational resources. Recent kd-tree-based approximate nearest neighbor (ANN) methods despite the sparsity of their compatibility tensor still require exhaustive calculations for large-scale graph matching. This work utilizes CUR tensor decomposition and introduces a novel cascaded second and third-order hypergraph matching framework (CURSOR) for efficient hypergraph matching. A CUR-based second-order graph matching algorithm is used to provide a rough match and then the core of CURSOR a fiber-CUR-based tensor generation method directly calculates entries of the compatibility tensor by leveraging the initial second-order match result. This significantly decreases the time complexity and tensor density. A probability relaxation labeling (PRL)-based matching algorithm especially suitable for sparse tensors is developed. Experiment results on large-scale synthetic datasets and widely-adopted benchmark sets demonstrate the superiority of CURSOR over existing methods. The tensor generation method in CURSOR can be integrated seamlessly into existing hypergraph matching methods to improve their performance and lower their computational costs.