CVPR 2024 论文列表

Machine learning models face generalization challenges when exposed to out-of-distribution (OOD) samples with unforeseen distribution shifts. Recent research reveals that for vision tasks test-time adaptation employing diffusion models can achieve state-of-the-art accuracy improvements on OOD samples by generating domain-aligned samples without altering the model's weights. Unfortunately those studies have primarily focused on pixel-level corruptions thereby lacking the generalization to adapt to a broader range of OOD types. We introduce Generalized Diffusion Adaptation (GDA) a novel diffusion-based test-time adaptation method robust against diverse OOD types. Specifically GDA iteratively guides the diffusion by applying a marginal entropy loss derived from the model in conjunction with style and content preservation losses during the reverse sampling process. In other words GDA considers the model's output behavior and the samples' semantic information as a whole reducing ambiguity in downstream tasks. based adaptation. Evaluation across various model architectures and OOD benchmarks indicates that GDA consistently surpasses previous diffusion-based adaptation methods. Notably it achieves the highest classification accuracy improvements ranging from 4.4% to 5.02% on ImageNet-C and 2.5% to 7.4% on Rendition Sketch and Stylized benchmarks. This performance highlights GDA's generalization to a broader range of OOD benchmarks.

We present GALA a framework that takes as input a single-layer clothed 3D human mesh and decomposes it into complete multi-layered 3D assets. The outputs can then be combined with other assets to create novel clothed human avatars with any pose. Existing reconstruction approaches often treat clothed humans as a single-layer of geometry and overlook the inherent compositionality of humans with hairstyles clothing and accessories thereby limiting the utility of the meshes for down-stream applications. Decomposing a single-layer mesh into separate layers is a challenging task because it requires the synthesis of plausible geometry and texture for the severely occluded regions. Moreover even with successful decomposition meshes are not normalized in terms of poses and body shapes failing coherent composition with novel identities and poses. To address these challenges we propose to leverage the general knowledge of a pretrained 2D diffusion model as geometry and appearance prior for humans and other assets. We first separate the input mesh using the 3D surface segmentation extracted from multi-view 2D segmentations. Then we synthesize the missing geometry of different layers in both posed and canonical spaces using a novel pose-guided Score Distillation Sampling (SDS) loss. Once we complete inpainting high-fidelity 3D geometry we also apply the same SDS loss to its texture to obtain the complete appearance including the initially occluded regions. Through a series of decomposition steps we obtain multiple layers of 3D assets in a shared canonical space normalized in terms of poses and human shapes hence supporting effortless composition to novel identities and reanimation with novel poses. Our experiments demonstrate the effectiveness of our approach for decomposition canonicalization and composition tasks compared to existing solutions.

We propose G-HOP a denoising diffusion based generative prior for hand-object interactions that allows modeling both the 3D object and a human hand conditioned on the object category. To learn a 3D spatial diffusion model that can capture this joint distribution we represent the human hand via a skeletal distance field to obtain a representation aligned with the (latent) signed distance field for the object. We show that this hand-object prior can then serve as a generic guidance to facilitate other tasks like reconstruction from interaction clip and human grasp synthesis. We believe that our model trained by aggregating several diverse real-world interaction datasets spanning 155 categories represents a first approach that allows jointly generating both hand and object. Our empirical evaluations demonstrate the benefit of this joint prior in video-based reconstruction and human grasp synthesis outperforming current task-specific baselines.

We propose functional diffusion a generative diffusion model focused on infinite-dimensional function data samples. In contrast to previous work functional diffusion works on samples that are represented by functions with a continuous domain. Functional diffusion can be seen as an extension of classical diffusion models to an infinite-dimensional domain. Functional diffusion is very versatile as images videos audio 3D shapes deformations etc. can be handled by the same framework with minimal changes. In addition functional diffusion is especially suited for irregular data or data defined in non-standard domains. In our work we derive the necessary foundations for functional diffusion and propose a first implementation based on the transformer architecture. We show generative results on complicated signed distance functions and deformation functions defined on 3D surfaces.

In this paper we uncover the untapped potential of diffusion U-Net which serves as a "free lunch" that substantially improves the generation quality on the fly. We initially investigate the key contributions of the U-Net architecture to the denoising process and identify that its main backbone primarily contributes to denoising whereas its skip connections mainly introduce high-frequency features into the decoder module causing the potential neglect of crucial functions intrinsic to the backbone network. Capitalizing on this discovery we propose a simple yet effective method termed "FreeU" which enhances generation quality without additional training or finetuning. Our key insight is to strategically re-weight the contributions sourced from the U-Net's skip connections and backbone feature maps to leverage the strengths of both components of the U-Net architecture. Promising results on image and video generation tasks demonstrate that our FreeU can be readily integrated to existing diffusion models e.g. Stable Diffusion DreamBooth and ControlNet to improve the generation quality with only a few lines of code. All you need is to adjust two scaling factors during inference.

Recent approaches such as ControlNet offer users fine-grained spatial control over text-to-image (T2I) diffusion models. However auxiliary modules have to be trained for each spatial condition type model architecture and checkpoint putting them at odds with the diverse intents and preferences a human designer would like to convey to the AI models during the content creation process. In this work we present FreeControl a training-free approach for controllable T2I generation that supports multiple conditions architectures and checkpoints simultaneously. FreeControl enforces structure guidance to facilitate the global alignment with a guidance image and appearance guidance to collect visual details from images generated without control. Extensive qualitative and quantitative experiments demonstrate the superior performance of FreeControl across a variety of pre-trained T2I models. In particular FreeControl enables convenient training-free control over many different architectures and checkpoints allows the challenging input conditions on which most of the existing training-free methods fail and achieves competitive synthesis quality compared to training-based approaches. Project page:https://genforce.github.io/freecontrol/.

Existing joint low-light enhancement and deblurring methods learn pixel-wise mappings from paired synthetic data which results in limited generalization in real-world scenes. While some studies explore the rich generative prior of pre-trained diffusion models they typically rely on the assumed degradation process and cannot handle unknown real-world degradations well. To address these problems we propose a novel zero-shot framework FourierDiff which embeds Fourier priors into a pre-trained diffusion model to harmoniously handle the joint degradation of luminance and structures. FourierDiff is appealing in its relaxed requirements on paired training data and degradation assumptions. The key zero-shot insight is motivated by image characteristics in the Fourier domain: most luminance information concentrates on amplitudes while structure and content information are closely related to phases. Based on this observation we decompose the sampled results of the reverse diffusion process in the Fourier domain and take advantage of the amplitude of the generative prior to align the enhanced brightness with the distribution of natural images. To yield a sharp and content-consistent enhanced result we further design a spatial-frequency alternating optimization strategy to progressively refine the phase of the input. Extensive experiments demonstrate the superior effectiveness of the proposed method especially in real-world scenes.

In this paper we study the problem of generalizable synthetic image detection aiming to detect forgery images from diverse generative methods e.g. GANs and diffusion models. Cutting-edge solutions start to explore the benefits of pre-trained models and mainly follow the fixed paradigm of solely training an attached classifier e.g. combining frozen CLIP-ViT with a learnable linear layer in UniFD. However our analysis shows that such a fixed paradigm is prone to yield detectors with insufficient learning regarding forgery representations. We attribute the key challenge to the lack of forgery adaptation and present a novel forgery-aware adaptive transformer approach namely FatFormer. Based on the pre-trained vision-language spaces of CLIP FatFormer introduces two core designs for the adaption to build generalized forgery representations. First motivated by the fact that both image and frequency analysis are essential for synthetic image detection we develop a forgery-aware adapter to adapt image features to discern and integrate local forgery traces within image and frequency domains. Second we find that considering the contrastive objectives between adapted image features and text prompt embeddings a previously overlooked aspect results in a nontrivial generalization improvement. Accordingly we introduce language-guided alignment to supervise the forgery adaptation with image and text prompts in FatFormer. Experiments show that by coupling these two designs our approach tuned on 4-class ProGAN data attains a remarkable detection performance achieving an average of 98% accuracy to unseen GANs and surprisingly generalizes to unseen diffusion models with 95% accuracy.

Diffusion models have transformed the image-to-image (I2I) synthesis and are now permeating into videos. However the advancement of video-to-video (V2V) synthesis has been hampered by the challenge of maintaining temporal consistency across video frames. This paper proposes a consistent V2V synthesis framework by jointly leveraging spatial conditions and temporal optical flow clues within the source video. Contrary to prior methods that strictly adhere to optical flow our approach harnesses its benefits while handling the imperfection in flow estimation. We encode the optical flow via warping from the first frame and serve it as a supplementary reference in the diffusion model. This enables our model for video synthesis by editing the first frame with any prevalent I2I models and then propagating edits to successive frames. Our V2V model FlowVid demonstrates remarkable properties: (1) Flexibility: FlowVid works seamlessly with existing I2I models facilitating various modifications including stylization object swaps and local edits. (2) Efficiency: Generation of a 4-second video with 30 FPS and 512x512 resolution takes only 1.5 minutes which is 3.1x 7.2x and 10.5x faster than CoDeF Rerender and TokenFlow respectively. (3) High-quality: In user studies our FlowVid is preferred 45.7% of the time outperforming CoDeF (3.5%) Rerender (10.2%) and TokenFlow (40.4%).

Optical flow estimation a process of predicting pixel-wise displacement between consecutive frames has commonly been approached as a regression task in the age of deep learning. Despite notable advancements this de facto paradigm unfortunately falls short in generalization performance when trained on synthetic or constrained data. Pioneering a paradigm shift we reformulate optical flow estimation as a conditional flow generation challenge unveiling FlowDiffuser --- a new family of optical flow models that could have stronger learning and generalization capabilities. FlowDiffuser estimates optical flow through a `noise-to-flow' strategy progressively eliminating noise from randomly generated flows conditioned on the provided pairs. To optimize accuracy and efficiency our FlowDiffuser incorporates a novel Conditional Recurrent Denoising Decoder (Conditional-RDD) streamlining the flow estimation process. It incorporates a unique Hidden State Denoising (HSD) paradigm effectively leveraging the information from previous time steps. Moreover FlowDiffuser can be easily integrated into existing flow networks leading to significant improvements in performance metrics compared to conventional implementations. Experiments on challenging benchmarks including Sintel and KITTI demonstrate the effectiveness of our FlowDiffuser with superior performance to existing state-of-the-art models. Code is available at https://github.com/LA30/FlowDiffuser.

In recent years there has been significant progress in the development of text-to-image generative models. Evaluating the quality of the generative models is one essential step in the development process. Unfortunately the evaluation process could consume a significant amount of computational resources making the required periodic evaluation of model performance (e.g. monitoring training progress) impractical. Therefore we seek to improve the evaluation efficiency by selecting the representative subset of the text-image dataset. We systematically investigate the design choices including the selection criteria (textural features or imagebased metrics) and the selection granularity (prompt-level or set-level). We find that the insights from prior work on subset selection for training data do not generalize to this problem and we propose FlashEval an iterative search algorithm tailored to evaluation data selection. We demonstrate the effectiveness of FlashEval on ranking diffusion models with various configurations including architectures quantization levels and sampler schedules on COCO and DiffusionDB datasets. Our searched 50-item subset could achieve comparable evaluation quality to the randomly sampled 500-item subset for COCO annotations on unseen models achieving a 10x evaluation speedup. We release the condensed subset of these commonly used datasets to help facilitate diffusion algorithm design and evaluation and open-source FlashEval as a tool for condensing future datasets accessible at https://github.com/thu-nics/FlashEval.

We introduce the Fixed Point Diffusion Model (FPDM) a novel approach to image generation that integrates the concept of fixed point solving into the framework of diffusion-based generative modeling. Our approach embeds an implicit fixed point solving layer into the denoising network of a diffusion model transforming the diffusion process into a sequence of closely-related fixed point problems. Combined with a new stochastic training method this approach significantly reduces model size reduces memory usage and accelerates training. Moreover it enables the development of two new techniques to improve sampling efficiency: reallocating computation across timesteps and reusing fixed point solutions between timesteps. We conduct extensive experiments with state-of-the-art models on ImageNet FFHQ CelebA-HQ and LSUN-Church demonstrating substantial improvements in performance and efficiency. Compared to the state-of-the-art DiT model FPDM contains 87% fewer parameters consumes 60% less memory during training and improves image generation quality in situations where sampling computation or time is limited.

The 3D Human Pose Estimation (3D HPE) task uses 2D images or videos to predict human joint coordinates in 3D space. Despite recent advancements in deep learning-based methods they mostly ignore the capability of coupling accessible texts and naturally feasible knowledge of humans missing out on valuable implicit supervision to guide the 3D HPE task. Moreover previous efforts often study this task from the perspective of the whole human body neglecting fine-grained guidance hidden in different body parts. To this end we present a new Fine-Grained Prompt-Driven Denoiser based on a diffusion model for 3D HPE named FinePOSE. It consists of three core blocks enhancing the reverse process of the diffusion model: (1) Fine-grained Part-aware Prompt learning (FPP) block constructs fine-grained part-aware prompts via coupling accessible texts and naturally feasible knowledge of body parts with learnable prompts to model implicit guidance. (2) Fine-grained Prompt-pose Communication (FPC) block establishes fine-grained communications between learned part-aware prompts and poses to improve the denoising quality. (3) Prompt-driven Timestamp Stylization (PTS) block integrates learned prompt embedding and temporal information related to the noise level to enable adaptive adjustment at each denoising step. Extensive experiments on public single-human pose estimation datasets show that FinePOSE outperforms state-of-the-art methods. We further extend FinePOSE to multi-human pose estimation. Achieving 34.3mm average MPJPE on the EgoHumans dataset demonstrates the potential of FinePOSE to deal with complex multi-human scenarios. Code is available at https://github.com/PKU-ICST-MIPL/FinePOSE_CVPR2024.

Even when using large multi-modal foundation models few-shot learning is still challenging -- if there is no proper inductive bias it is nearly impossible to keep the nuanced class attributes while removing the visually prominent attributes that spuriously correlate with class labels. To this end we find an inductive bias that the time-steps of a Diffusion Model (DM) can isolate the nuanced class attributes i.e. as the forward diffusion adds noise to an image at each time-step nuanced attributes are usually lost at an earlier time-step than the spurious attributes that are visually prominent. Building on this we propose Time-step Few-shot (TiF) learner. We train class-specific low-rank adapters for a text-conditioned DM to make up for the lost attributes such that images can be accurately reconstructed from their noisy ones given a prompt. Hence at a small time-step the adapter and prompt are essentially a parameterization of only the nuanced class attributes. For a test image we can use the parameterization to only extract the nuanced class attributes for classification. TiF learner significantly outperforms OpenCLIP and its adapters on a variety of fine-grained and customized few-shot learning tasks. Codes are in https://github.com/yue-zhongqi/tif.

Sampling from diffusion models can be treated as solving the corresponding ordinary differential equations (ODEs) with the aim of obtaining an accurate solution with as few number of function evaluations (NFE) as possible. Recently various fast samplers utilizing higher-order ODE solvers have emerged and achieved better performance than the initial first-order one. However these numerical methods inherently result in certain approximation errors which significantly degrades sample quality with extremely small NFE (e.g. around 5). In contrast based on the geometric observation that each sampling trajectory almost lies in a two-dimensional subspace embedded in the ambient space we propose Approximate MEan-Direction Solver (AMED-Solver) that eliminates truncation errors by directly learning the mean direction for fast diffusion sampling. Besides our method can be easily used as a plugin to further improve existing ODE-based samplers. Extensive experiments on image synthesis with the resolution ranging from 32 to 512 demonstrate the effectiveness of our method. With only 5 NFE we achieve 6.61 FID on CIFAR-10 10.74 FID on ImageNet 64x64 and 13.20 FID on LSUN Bedroom. Our code is available at https://github.com/zju-pi/diff-sampler.

Due to the high potential for abuse of GenAI systems the task of detecting synthetic images has recently become of great interest to the research community. Unfortunately existing image space detectors quickly become obsolete as new high-fidelity text-to-image models are developed at blinding speed. In this work we propose a new synthetic image detector that uses features obtained by inverting an open-source pre-trained Stable Diffusion model. We show that these inversion features enable our detector to generalize well to unseen generators of high visual fidelity (e.g. DALL*E 3) even when the detector is trained only on lower fidelity fake images generated via Stable Diffusion. This detector achieves new state-of-the-art across multiple training and evaluation setups. Moreover we introduce a new challenging evaluation protocol that uses reverse image search to mitigate stylistic and thematic biases in the detector evaluation. We show that the resulting evaluation scores align well with detectors' in-the-wild performance and release these datasets as public benchmarks for future research.

In this paper we introduce Fairy a minimalist yet robust adaptation of image-editing diffusion models enhancing them for video editing applications. Our approach centers on the concept of anchor-based cross-frame attention a mechanism that implicitly propagates diffusion features across frames ensuring superior temporal coherence and high-fidelity synthesis. Fairy not only addresses limitations of previous models including memory and processing speed. It also improves temporal consistency through a unique data augmentation strategy. This strategy renders the model equivariant to affine transformations in both source and target images. Remarkably efficient Fairy generates 120-frame 512x384 videos (4-second duration at 30 FPS) in just 14 seconds outpacing prior works by at least 44x. A comprehensive user study involving 1000 generated samples confirms that our approach delivers superior quality decisively outperforming established methods.

The goal of this work is to simultaneously generate natural talking faces and speech outputs from text. We achieve this by integrating Talking Face Generation (TFG) and Text-to-Speech (TTS) systems into a unified framework. We address the main challenges of each task: (1) generating a range of head poses representative of real-world scenarios and (2) ensuring voice consistency despite variations in facial motion for the same identity. To tackle these issues we introduce a motion sampler based on conditional flow matching which is capable of high-quality motion code generation in an efficient way. Moreover we introduce a novel conditioning method for the TTS system which utilises motion-removed features from the TFG model to yield uniform speech outputs. Our extensive experiments demonstrate that our method effectively creates natural-looking talking faces and speech that accurately match the input text. To our knowledge this is the first effort to build a multimodal synthesis system that can generalise to unseen identities.

We introduce FaceTalk a novel generative approach designed for synthesizing high-fidelity 3D motion sequences of talking human heads from input audio signal. To capture the expressive detailed nature of human heads including hair ears and finer-scale eye movements we propose to couple speech signal with the latent space of neural parametric head models to create high-fidelity temporally coherent motion sequences. We propose a new latent diffusion model for this task operating in the expression space of neural parametric head models to synthesize audio-driven realistic head sequences. In the absence of a dataset with corresponding NPHM expressions to audio we optimize for these correspondences to produce a dataset of temporally-optimized NPHM expressions fit to audio-video recordings of people talking. To the best of our knowledge this is the first work to propose a generative approach for realistic and high-quality motion synthesis of volumetric human heads representing a significant advancement in the field of audio-driven 3D animation. Notably our approach stands out in its ability to generate plausible motion sequences that can produce high-fidelity head animation coupled with the NPHM shape space. Our experimental results substantiate the effectiveness of FaceTalk consistently achieving superior and visually natural motion encompassing diverse facial expressions and styles outperforming existing methods by 75% in perceptual user study evaluation