扩散模型

We introduce pix2gestalt a framework for zero-shot amodal segmentation which learns to estimate the shape and appearance of whole objects that are only partially visible behind occlusions. By capitalizing on large-scale diffusion models and transferring their representations to this task we learn a conditional diffusion model for reconstructing whole objects in challenging zero-shot cases including examples that break natural and physical priors such as art. As training data we use a synthetically curated dataset containing occluded objects paired with their whole counterparts. Experiments show that our approach outperforms supervised baselines on established benchmarks. Our model can furthermore be used to significantly improve the performance of existing object recognition and 3D reconstruction methods in the presence of occlusions.

iToF is a prevalent cost-effective technology for 3D perception. While its reliance on multi-measurement commonly leads to reduced performance in dynamic environments. Based on the analysis of the physical iToF imaging process we propose the iToF flow composed of crossmode transformation and uni-mode photometric correction to model the variation of measurements caused by different measurement modes and 3D motion respectively. We propose a local linear transform (LLT) based cross-mode transfer module (LCTM) for mode-varying and pixel shift compensation of cross-mode flow and uni-mode photometric correct module (UPCM) for estimating the depth-wise motion caused photometric residual of uni-mode flow. The iToF flow-based depth extraction network is proposed which could facilitate the estimation of the 4-phase measurements at each individual time for high framerate and accurate depth estimation. Extensive experiments including both simulation and real-world experiments are conducted to demonstrate the effectiveness of the proposed methods. Compared with the SOTA method our approach reduces the computation time by 75% while improving the performance by 38%. The code and database are available at https://github.com/ComputationalPerceptionLab/iToF_flow.

We introduce a 3D-aware diffusion model ZeroNVS for single-image novel view synthesis for in-the-wild scenes. While existing methods are designed for single objects with masked backgrounds we propose new techniques to address challenges introduced by in-the-wild multi-object scenes with complex backgrounds. Specifically we train a generative prior on a mixture of data sources that capture object-centric indoor and outdoor scenes. To address issues from data mixture such as depth-scale ambiguity we propose a novel camera conditioning parameterization and normalization scheme. Further we observe that Score Distillation Sampling (SDS) tends to truncate the distribution of complex backgrounds during distillation of 360-degree scenes and propose "SDS anchoring" to improve the diversity of synthesized novel views. Our model sets a new state-of-the-art result in LPIPS on the DTU dataset in the zero-shot setting even outperforming methods specifically trained on DTU. We further adapt the challenging Mip-NeRF 360 dataset as a new benchmark for single-image novel view synthesis and demonstrate strong performance in this setting. Code and models will be publicly available.

Tone mapping techniques aiming to convert high dynamic range (HDR) images to high-quality low dynamic range (LDR) images for display play a more crucial role in real-world vision systems with the increasing application of HDR images. However obtaining paired HDR and high-quality LDR images is difficult posing a challenge to deep learning based tone mapping methods. To overcome this challenge we propose a novel zero-shot tone mapping framework that utilizes shared structure knowledge allowing us to transfer a pre-trained mapping model from the LDR domain to HDR fields without paired training data. Our approach involves decomposing both the LDR and HDR images into two components: structural information and tonal information. To preserve the original image's structure we modify the reverse sampling process of a diffusion model and explicitly incorporate the structure information into the intermediate results. Additionally for improved image details we introduce a dual-control network architecture that enables different types of conditional inputs to control different scales of the output. Experimental results demonstrate the effectiveness of our approach surpassing previous state-of-the-art methods both qualitatively and quantitatively. Moreover our model exhibits versatility and can be applied to other low-level vision tasks without retraining. The code is available at https://github.com/ZSDM-HDR/Zero-Shot-Diffusion-HDR.

Understanding illumination and reducing the need for supervision pose a significant challenge in low-light enhancement. Current approaches are highly sensitive to data usage during training and illumination-specific hyper-parameters limiting their ability to handle unseen scenarios. In this paper we propose a new zero-reference low-light enhancement framework trainable solely with normal light images. To accomplish this we devise an illumination-invariant prior inspired by the theory of physical light transfer. This prior serves as the bridge between normal and low-light images. Then we develop a prior-to-image framework trained without low-light data. During testing this framework is able to restore our illumination-invariant prior back to images automatically achieving low-light enhancement. Within this framework we leverage a pretrained generative diffusion model for model ability introduce a bypass decoder to handle detail distortion as well as offer a lightweight version for practicality. Extensive experiments demonstrate our framework's superiority in various scenarios as well as good interpretability robustness and efficiency. Code is available on our project homepage: http://daooshee.github.io/QuadPrior-Website/

Despite the remarkable progress in image style transfer formulating style in the context of art is inherently subjective and challenging. In contrast to existing methods this study shows that vanilla diffusion models can directly extract style information and seamlessly integrate the generative prior into the content image without retraining. Specifically we adopt dual denoising paths to represent content/style references in latent space and then guide the content image denoising process with style latent codes. We further reveal that the cross-attention mechanism in latent diffusion models tends to blend the content and style images resulting in stylized outputs that deviate from the original content image. To overcome this limitation we introduce a cross-attention reweighting strategy. Through theoretical analysis and experiments we demonstrate the effectiveness and superiority of the diffusion-based zero-shot style transfer via attention reweighting Z-STAR.

Knowledge distillation methods have recently shown to be a promising direction to speedup the synthesis of large-scale diffusion models by requiring only a few inference steps. While several powerful distillation methods were recently proposed the overall quality of student samples is typically lower compared to the teacher ones which hinders their practical usage. In this work we investigate the relative quality of samples produced by the teacher text-to-image diffusion model and its distilled student version. As our main empirical finding we discover that a noticeable portion of student samples exhibit superior fidelity compared to the teacher ones despite the approximate nature of the student. Based on this finding we propose an adaptive collaboration between student and teacher diffusion models for effective text-to-image synthesis. Specifically the distilled model produces an initial image sample and then an oracle decides whether it needs further improvements with the teacher model. Extensive experiments demonstrate that the designed pipeline surpasses state-of-the-art text-to-image alternatives for various inference budgets in terms of human preference. Furthermore the proposed approach can be naturally used in popular applications such as text-guided image editing and controllable generation.

We present XCube a novel generative model for high-resolution sparse 3D voxel grids with arbitrary attributes. Our model can generate millions of voxels with a finest effective resolution of up to 1024^3 in a feed-forward fashion without time-consuming test-time optimization. To achieve this we employ a hierarchical voxel latent diffusion model which generates progressively higher resolution grids in a coarse-to-fine manner using a custom framework built on the highly efficient VDB data structure. Apart from generating high-resolution objects we demonstrate the effectiveness of XCube on large outdoor scenes at scales of 100m x 100m with a voxel size as small as 10cm. We observe clear qualitative and quantitative improvements over past approaches. In addition to unconditional generation we show that our model can be used to solve a variety of tasks such as user-guided editing scene completion from a single scan and text-to-3D.

We introduce X-Adapter a universal upgrader to enable the pretrained plug-and-play modules (e.g. ControlNet LoRA) to work directly with the upgraded text-to-image diffusion model (e.g. SDXL) without further retraining. We achieve this goal by training an additional network to control the frozen upgraded model with the new text-image data pairs. In detail X-Adapter keeps a frozen copy of the old model to preserve the connectors of different plugins. Additionally X-Adapter adds trainable mapping layers that bridge the decoders from models of different versions for feature remapping. The remapped features will be used as guidance for the upgraded model. To enhance the guidance ability of X-Adapter we employ a -text training strategy for the upgraded model. After training we also introduce a two-stage denoising strategy to align the initial latents of X-Adapter and the upgraded model. Thanks to our strategies X-Adapter demonstrates universal compatibility with various plugins and also enables plugins of different versions to work together thereby expanding the functionalities of diffusion community. To verify the effectiveness of the proposed method we conduct extensive experiments and the results show that X-Adapter may facilitate wider application in the upgraded foundational diffusion model. Project page at: https://showlab.github.io/X-Adapter.

In this work we introduce Wonder3D a novel method for generating high-fidelity textured meshes from single-view images with remarkable efficiency. Recent methods based on the Score Distillation Sampling (SDS) loss methods have shown the potential to recover 3D geometry from 2D diffusion priors but they typically suffer from time-consuming per-shape optimization and inconsistent geometry. In contrast certain works directly produce 3D information via fast network inferences but their results are often of low quality and lack geometric details. To holistically improve the quality consistency and efficiency of image-to-3D tasks we propose a cross-domain diffusion model that generates multi-view normal maps and the corresponding color images. To ensure consistency we employ a multi-view cross-domain attention mechanism that facilitates information exchange across views and modalities. Lastly we introduce a geometry-aware normal fusion algorithm that extracts high-quality surfaces from the multi-view 2D representations in only 2 3 minutes. Our extensive evaluations demonstrate that our method achieves high-quality reconstruction results robust generalization and remarkable efficiency compared to prior works.

Creating multi-view wire art (MVWA) a static 3D sculpture with diverse interpretations from different viewpoints is a complex task even for skilled artists. In response we present DreamWire an AI system enabling everyone to craft MVWA easily. Users express their vision through text prompts or scribbles freeing them from intricate 3D wire organisation. Our approach synergises 3D Bezier curves Prim's algorithm and knowledge distillation from diffusion models or their variants (e.g. ControlNet). This blend enables the system to represent 3D wire art ensuring spatial continuity and overcoming data scarcity. Extensive evaluation and analysis are conducted to shed insight on the inner workings of the proposed system including the trade-off between connectivity and visual aesthetics.

Text-to-image diffusion models have remarkably excelled in producing diverse high-quality and photo-realistic images. This advancement has spurred a growing interest in incorporating specific identities into generated content. Most current methods employ an inversion approach to embed a target visual concept into the text embedding space using a single reference image. However the newly synthesized faces either closely resemble the reference image in terms of facial attributes such as expression or exhibit a reduced capacity for identity preservation. Text descriptions intended to guide the facial attributes of the synthesized face may fall short owing to the intricate entanglement of identity information with identity-irrelevant facial attributes derived from the reference image. To address these issues we present the novel use of the extended StyleGAN embedding space \mathcal W _+ to achieve enhanced identity preservation and disentanglement for diffusion models. By aligning this semantically meaningful human face latent space with text-to-image diffusion models we succeed in maintaining high fidelity in identity preservation coupled with the capacity for semantic editing. Additionally we propose new training objectives to balance the influences of both prompt and identity conditions ensuring that the identity-irrelevant background remains \lxm negligibly affected during facial attribute modifications. Extensive experiments reveal that our method adeptly generates personalized text-to-image outputs that are not only compatible with prompt descriptions but also amenable to common StyleGAN editing directions in diverse settings. Our code and model are available at https://github.com/csxmli2016/w-plus-adapter.

Underwater images are subject to intricate and diverse degradation inevitably affecting the effectiveness of underwater visual tasks. However most approaches primarily operate in the raw pixel space of images which limits the exploration of the frequency characteristics of underwater images leading to an inadequate utilization of deep models' representational capabilities in producing high-quality images. In this paper we introduce a novel Underwater Image Enhancement (UIE) framework named WF-Diff designed to fully leverage the characteristics of frequency domain information and diffusion models. WF-Diff consists of two detachable networks: Wavelet-based Fourier information interaction network (WFI2-net) and Frequency Residual Diffusion Adjustment Module (FRDAM). With our full exploration of the frequency domain information WFI2-net aims to achieve preliminary enhancement of frequency information in the wavelet space. Our proposed FRDAM can further refine the high- and low-frequency information of the initial enhanced images which can be viewed as a plug-and-play universal module to adjust the detail of the underwater images. With the above techniques our algorithm can show SOTA performance on real-world underwater image datasets and achieves competitive performance in visual quality.

Although diffusion models are rising as a powerful solution for blind face restoration they are criticized for two problems: 1) slow training and inference speed and 2) failure in preserving identity and recovering fine-grained facial details. In this work we propose WaveFace to solve the problems in the frequency domain where low- and high-frequency components decomposed by wavelet transformation are considered individually to maximize authenticity as well as efficiency. The diffusion model is applied to recover the low-frequency component only which presents general information of the original image but 1/16 in size. To preserve the original identity the generation is conditioned on the low-frequency component of low-quality images at each denoising step. Meanwhile high-frequency components at multiple decomposition levels are handled by a unified network which recovers complex facial details in a single step. Evaluations on four benchmark datasets show that: 1) WaveFace outperforms state-of-the-art methods in authenticity especially in terms of identity preservation and 2) authentic images are restored with the efficiency 10x faster than existing diffusion model-based BFR methods.

Diffusion Models (DMs) have shown remarkable capabilities in various image-generation tasks. However there are growing concerns that DMs could be used to imitate unauthorized creations and thus raise copyright issues. To address this issue we propose a novel framework that embeds personal watermarks in the generation of adversarial examples. Such examples can force DMs to generate images with visible watermarks and prevent DMs from imitating unauthorized images. We construct a generator based on conditional adversarial networks and design three losses (adversarial loss GAN loss and perturbation loss) to generate adversarial examples that have subtle perturbation but can effectively attack DMs to prevent copyright violations. Training a generator for a personal watermark by our method only requires 5-10 samples within 2-3 minutes and once the generator is trained it can generate adversarial examples with that watermark significantly fast (0.2s per image). We conduct extensive experiments in various conditional image-generation scenarios. Compared to existing methods that generate images with chaotic textures our method adds visible watermarks on the generated images which is a more straightforward way to indicate copyright violations. We also observe that our adversarial examples exhibit good transferability across unknown generative models. Therefore this work provides a simple yet powerful way to protect copyright from DM-based imitation.

The rapid advancement of generative models facilitating the creation of hyper-realistic images from textual descriptions has concurrently escalated critical societal concerns such as misinformation. Although providing some mitigation traditional fingerprinting mechanisms fall short in attributing responsibility for the malicious use of synthetic images. This paper introduces a novel approach to model fingerprinting that assigns responsibility for the generated images thereby serving as a potential countermeasure to model misuse. Our method modifies generative models based on each user's unique digital fingerprint imprinting a unique identifier onto the resultant content that can be traced back to the user. This approach incorporating fine-tuning into Text-to-Image (T2I) tasks using the Stable Diffusion Model demonstrates near-perfect attribution accuracy with a minimal impact on output quality. Through extensive evaluation we show that our method outperforms baseline methods with an average improvement of 11% in handling image post-processes. Our method presents a promising and novel avenue for accountable model distribution and responsible use. Our code is available in https://github.com/kylemin/WOUAF.

In this work we present Vlogger a generic AI system for generating a minute-level video blog (i.e. vlog) of user descriptions. Different from short videos with a few seconds vlog often contains a complex storyline with diversified scenes which is challenging for most existing video generation approaches. To break through this bottleneck our Vlogger smartly leverages Large Language Model (LLM) as Director and decomposes a long video generation task of vlog into four key stages where we invoke various foundation models to play the critical roles of vlog professionals including (1) Script (2) Actor (3) ShowMaker and (4) Voicer. With such a design of mimicking human beings our Vlogger can generate vlogs through explainable cooperation of top-down planning and bottom-up shooting. More over we introduce a novel video diffusion model ShowMaker which serves as a videographer in our Vlogger for generating the video snippet of each shooting scene. By incorporating Script and Actor attentively as textual and visual prompts it can effectively enhance spatial-temporal coherence in the snippet. Besides we design a concise mixed training paradigm for ShowMaker boosting its capacity for both T2V generation and prediction. Finally the extensive experiments show that our method achieves state-of-the-art performance on zero-shot T2V generation and prediction tasks. More importantly Vlogger can generate over 5-minute vlogs from open-world descriptions without loss of video coherence on script and actor.

This paper proposes an image-vector dual diffusion model for generative layout design. Distinct from prior efforts that mostly ignore element-level visual information our approach integrates the power of a pre-trained large image diffusion model to guide layout composition in a vector diffusion model by providing enhanced salient region understanding and high-level inter-element relationship reasoning. Our proposed model simultaneously operates in two domains: it generates the overall design appearance in the image domain while optimizing the size and position of each design element in the vector domain. The proposed method achieves the state-of-the-art results on several datasets and enables new layout design applications.

We address the problem of synthesizing multi-view optical illusions: images that change appearance upon a transformation such as a flip or rotation. We propose a simple zero-shot method for obtaining these illusions from off-the-shelf text-to-image diffusion models. During the reverse diffusion process we estimate the noise from different views of a noisy image and then combine these noise estimates together and denoise the image. A theoretical analysis suggests that this method works precisely for views that can be written as orthogonal transformations of which permutations are a subset. This leads to the idea of a visual anagram ---an image that changes appearance under some rearrangement of pixels. This includes rotations and flips but also more exotic pixel permutations such as a jigsaw rearrangement. Our approach also naturally extends to illusions with more than two views. We provide both qualitative and quantitative results demonstrating the effectiveness and flexibility of our method. Please see our project webpage for additional visualizations and results: https://dangeng.github.io/visual_anagrams/

Novel-view synthesis through diffusion models has demonstrated remarkable potential for generating diverse and high-quality images. Yet the independent process of image generation in these prevailing methods leads to challenges in maintaining multiple-view consistency. To address this we introduce ViewFusion a novel training-free algorithm that can be seamlessly integrated into existing pre-trained diffusion models. Our approach adopts an auto-regressive method that implicitly leverages previously generated views as context for the next view generation ensuring robust multi-view consistency during the novel-view generation process. Through a diffusion process that fuses known-view information via interpolated denoising our framework successfully extends single-view conditioned models to work in multiple-view conditional settings without any additional fine-tuning. Extensive experimental results demonstrate the effectiveness of ViewFusion in generating consistent and detailed novel views.