CVPR 2024 Papers

Conference on Computer Vision and Pattern Recognition

Decision breakdown

2,716papers

Counted from the official paper list; may differ slightly from the number the conference announced.

  • Poster
    2,39288.1%
  • Other / unlabeled
    32411.9%

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Topic: Diffusion Models · 496 papers · Reset filters

Style Injection in Diffusion: A Training-free Approach for Adapting Large-scale Diffusion Models for Style Transfer

Jiwoo Chung, Sangeek Hyun, Jae-Pil Heo

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Despite the impressive generative capabilities of diffusion models existing diffusion model-based style transfer methods require inference-stage optimization (e.g. fine-tuning or textual inversion of style) which is time-consuming or fails to leverage the generative ability of large-scale diffusion models. To address these issues we introduce a novel artistic style transfer method based on a pre-trained large-scale diffusion model without any optimization. Specifically we manipulate the features of self-attention layers as the way the cross-attention mechanism works; in the generation process substituting the key and value of content with those of style image. This approach provides several desirable characteristics for style transfer including 1) preservation of content by transferring similar styles into similar image patches and 2) transfer of style based on similarity of local texture (e.g. edge) between content and style images. Furthermore we introduce query preservation and attention temperature scaling to mitigate the issue of disruption of original content and initial latent Adaptive Instance Normalization (AdaIN) to deal with the disharmonious color (failure to transfer the colors of style). Our experimental results demonstrate that our proposed method surpasses state-of-the-art methods in both conventional and diffusion-based style transfer baselines.

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Structure-Guided Adversarial Training of Diffusion Models

Ling Yang, Haotian Qian, Zhilong Zhang, Jingwei Liu, Bin Cui

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Diffusion models have demonstrated exceptional efficacy in various generative applications. While existing models focus on minimizing a weighted sum of denoising score matching losses for data distribution modeling their training primarily emphasizes instance-level optimization overlooking valuable structural information within each mini-batch indicative of pair-wise relationships among samples. To address this limitation we introduce Structure-guided Adversarial training of Diffusion Models (SADM). In this pioneering approach we compel the model to learn manifold structures between samples in each training batch. To ensure the model captures authentic manifold structures in the data distribution we advocate adversarial training of the diffusion generator against a novel structure discriminator in a minimax game distinguishing real manifold structures from the generated ones. SADM substantially outperforms existing methods in image generation and cross-domain fine-tuning tasks across 12 datasets establishing a new state-of-the-art FID of 1.58 and 2.11 on ImageNet for class-conditional image generation at resolutions of 256x256 and 512x512 respectively.

Poster
Stratified Avatar Generation from Sparse Observations

Han Feng, Wenchao Ma, Quankai Gao, Xianwei Zheng, Nan Xue, Huijuan Xu

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Estimating 3D full-body avatars from AR/VR devices is essential for creating immersive experiences in AR/VR applications. This task is challenging due to the limited input from Head Mounted Devices which capture only sparse observations from the head and hands. Predicting the full-body avatars particularly the lower body from these sparse observations presents significant difficulties. In this paper we are inspired by the inherent property of the kinematic tree defined in the Skinned Multi-Person Linear (SMPL) model where the upper body and lower body share only one common ancestor node bringing the potential of decoupled reconstruction. We propose a stratified approach to decouple the conventional full-body avatar reconstruction pipeline into two stages with the reconstruction of the upper body first and a subsequent reconstruction of the lower body conditioned on the previous stage. To implement this straightforward idea we leverage the latent diffusion model as a powerful probabilistic generator and train it to follow the latent distribution of decoupled motions explored by a VQ-VAE encoder-decoder model. Extensive experiments on AMASS mocap dataset demonstrate our state-of-the-art performance in the reconstruction of full-body motions.

Poster
Space-Time Diffusion Features for Zero-Shot Text-Driven Motion Transfer

Danah Yatim, Rafail Fridman, Omer Bar-Tal, Yoni Kasten, Tali Dekel

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We present a new method for text-driven motion transfer - synthesizing a video that complies with an input text prompt describing the target objects and scene while maintaining an input video's motion and scene layout. Prior methods are confined to transferring motion across two subjects within the same or closely related object categories and are applicable for limited domains (e.g. humans). In this work we consider a significantly more challenging setting in which the target and source objects differ drastically in shape and fine-grained motion characteristics (e.g. translating a jumping dog into a dolphin). To this end we leverage a pre-trained and fixed text-to-video diffusion model which provides us with generative and motion priors. The pillar of our method is a new space-time feature loss derived directly from the model. This loss guides the generation process to preserve the overall motion of the input video while complying with the target object in terms of shape and fine-grained motion traits.

Poster
SonicVisionLM: Playing Sound with Vision Language Models

Zhifeng Xie, Shengye Yu, Qile He, Mengtian Li

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There has been a growing interest in the task of generating sound for silent videos primarily because of its practicality in streamlining video post-production. However existing methods for video-sound generation attempt to directly create sound from visual representations which can be challenging due to the difficulty of aligning visual representations with audio representations. In this paper we present SonicVisionLM a novel framework aimed at generating a wide range of sound effects by leveraging vision-language models(VLMs). Instead of generating audio directly from video we use the capabilities of powerful VLMs. When provided with a silent video our approach first identifies events within the video using a VLM to suggest possible sounds that match the video content. This shift in approach transforms the challenging task of aligning image and audio into more well-studied sub-problems of aligning image-to-text and text-to-audio through the popular diffusion models. To improve the quality of audio recommendations with LLMs we have collected an extensive dataset that maps text descriptions to specific sound effects and developed a time-controlled audio adapter. Our approach surpasses current state-of-the-art methods for converting video to audio enhancing synchronization with the visuals and improving alignment between audio and video components. Project page: https://yusiissy.github.io/SonicVisionLM.github.io/

Poster
Solving Masked Jigsaw Puzzles with Diffusion Vision Transformers

Jinyang Liu, Wondmgezahu Teshome, Sandesh Ghimire, Mario Sznaier, Octavia Camps

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Solving image and video jigsaw puzzles poses the challenging task of rearranging image fragments or video frames from unordered sequences to restore meaningful images and video sequences. Existing approaches often hinge on discriminative models tasked with predicting either the absolute positions of puzzle elements or the permutation actions applied to the original data. Unfortunately these methods face limitations in effectively solving puzzles with a large number of elements. In this paper we propose JPDVT an innovative approach that harnesses diffusion transformers to address this challenge. Specifically we generate positional information for image patches or video frames conditioned on their underlying visual content. This information is then employed to accurately assemble the puzzle pieces in their correct positions even in scenarios involving missing pieces. Our method achieves state-of-the-art performance on several datasets.

Poster
SocialCounterfactuals: Probing and Mitigating Intersectional Social Biases in Vision-Language Models with Counterfactual Examples

Phillip Howard, Avinash Madasu, Tiep Le, Gustavo Lujan Moreno, Anahita Bhiwandiwalla, Vasudev Lal

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While vision-language models (VLMs) have achieved remarkable performance improvements recently there is growing evidence that these models also posses harmful biases with respect to social attributes such as gender and race. Prior studies have primarily focused on probing such bias attributes individually while ignoring biases associated with intersections between social attributes. This could be due to the difficulty of collecting an exhaustive set of image-text pairs for various combinations of social attributes. To address this challenge we employ text-to-image diffusion models to produce counterfactual examples for probing intersectional social biases at scale. Our approach utilizes Stable Diffusion with cross attention control to produce sets of counterfactual image-text pairs that are highly similar in their depiction of a subject (e.g. a given occupation) while differing only in their depiction of intersectional social attributes (e.g. race & gender). Through our over-generate-then-filter methodology we produce SocialCounterfactuals a high-quality dataset containing 171k image-text pairs for probing intersectional biases related to gender race and physical characteristics. We conduct extensive experiments to demonstrate the usefulness of our generated dataset for probing and mitigating intersectional social biases in state-of-the-art VLMs.

Poster
Smooth Diffusion: Crafting Smooth Latent Spaces in Diffusion Models

Jiayi Guo, Xingqian Xu, Yifan Pu, Zanlin Ni, Chaofei Wang, Manushree Vasu, Shiji Song, Gao Huang, Humphrey Shi

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Recently diffusion models have made remarkable progress in text-to-image (T2I) generation synthesizing images with high fidelity and diverse contents. Despite this advancement latent space smoothness within diffusion models remains largely unexplored. Smooth latent spaces ensure that a perturbation on an input latent corresponds to a steady change in the output image. This property proves beneficial in downstream tasks including image interpolation inversion and editing. In this work we expose the non-smoothness of diffusion latent spaces by observing noticeable visual fluctuations resulting from minor latent variations. To tackle this issue we propose Smooth Diffusion a new category of diffusion models that can be simultaneously high-performing and smooth. Specifically we introduce Step-wise Variation Regularization to enforce the proportion between the variations of an arbitrary input latent and that of the output image is a constant at any diffusion training step. In addition we devise an interpolation standard deviation (ISTD) metric to effectively assess the latent space smoothness of a diffusion model. Extensive quantitative and qualitative experiments demonstrate that Smooth Diffusion stands out as a more desirable solution not only in T2I generation but also across various downstream tasks. Smooth Diffusion is implemented as a plug-and-play Smooth-LoRA to work with various community models. Code is available at https://github.com/SHI-Labs/Smooth-Diffusion.

Poster
SmartMask: Context Aware High-Fidelity Mask Generation for Fine-grained Object Insertion and Layout Control

Jaskirat Singh, Jianming Zhang, Qing Liu, Cameron Smith, Zhe Lin, Liang Zheng

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The field of generative image inpainting and object insertion has made significant progress with the recent advent of latent diffusion models. Utilizing a precise object mask can greatly enhance these applications. However due to the challenges users encounter in creating high-fidelity masks there is a tendency for these methods to rely on more coarse masks (e.g. bounding box) for these applications. This results in limited control and compromised background content preservation. To overcome these limitations we introduce SmartMask which allows any novice user to create detailed masks for precise object insertion. Combined with a ControlNet-Inpaint model our experiments demonstrate that SmartMask achieves superior object insertion quality preserving the background content more effectively than previous methods. Notably unlike prior works the proposed approach can also be used even without user-mask guidance which allows it to perform mask-free object insertion at diverse positions and scales. Furthermore we find that when used iteratively with a novel instruction-tuning based planning model SmartMask can be used to design detailed layouts from scratch. As compared with user-scribble based layout design we observe that SmartMask allows for better quality outputs with layout-to-image generation methods.

Poster
SmartEdit: Exploring Complex Instruction-based Image Editing with Multimodal Large Language Models

Yuzhou Huang, Liangbin Xie, Xintao Wang, Ziyang Yuan, Xiaodong Cun, Yixiao Ge, Jiantao Zhou, Chao Dong, Rui Huang, Ruimao Zhang, Ying Shan

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Current instruction-based image editing methods such as InstructPix2Pix often fail to produce satisfactory results in complex scenarios due to their dependence on the simple CLIP text encoder in diffusion models. To rectify this this paper introduces SmartEdit a novel approach of instruction-based image editing that leverages Multimodal Large Language Models (MLLMs) to enhance its understanding and reasoning capabilities. However direct integration of these elements still faces challenges in situations requiring complex reasoning. To mitigate this we propose a Bidirectional Interaction Module (BIM) that enables comprehensive bidirectional information interactions between the input image and the MLLM output. During training we initially incorporate perception data to boost the perception and understanding capabilities of diffusion models. Subsequently we demonstrate that a small amount of complex instruction editing data can effectively stimulate SmartEdit's editing capabilities for more complex instructions. We further construct a new evaluation dataset Reason-Edit specifically tailored for complex instruction-based image editing. Both quantitative and qualitative results on this evaluation dataset indicate that our SmartEdit surpasses previous methods paving the way for the practical application of complex instruction-based image editing.

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Slice3D: Multi-Slice Occlusion-Revealing Single View 3D Reconstruction

Yizhi Wang, Wallace Lira, Wenqi Wang, Ali Mahdavi-Amiri, Hao Zhang

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We introduce multi-slice reasoning a new notion for single-view 3D reconstruction which challenges the current and prevailing belief that multi-view synthesis is the most natural conduit between single-view and 3D. Our key observation is that object slicing is a more direct and hence more advantageous means to reveal occluded structures than altering camera views. Specifically slicing can peel through any occluder without obstruction and in the limit (i.e. with infinitely many slices) it is guaranteed to unveil all hidden object parts. We realize our idea by developing Slice3D a novel method for single-view 3D reconstruction which first predicts multi-slice images from a single RGB input image and then integrates the slices into a 3D model using a coordinate-based transformer network to product a signed distance function. The slice images can be regressed or generated both through a U-Net based network. For the former we inject a learnable slice indicator code to designate each decoded image into a spatial slice location while the slice generator is a denoising diffusion model operating on the entirety of slice images stacked on the input channels. We conduct extensive evaluation against state-of-the-art alternatives to demonstrate superiority of our method especially in recovering complex and severely occluded shape structures amid ambiguities. All Slice3D results were produced by networks trained on a single Nvidia A40 GPU with an inference time of less than 20 seconds.

Poster
SkillDiffuser: Interpretable Hierarchical Planning via Skill Abstractions in Diffusion-Based Task Execution

Zhixuan Liang, Yao Mu, Hengbo Ma, Masayoshi Tomizuka, Mingyu Ding, Ping Luo

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Diffusion models have demonstrated strong potential for robotic trajectory planning. However generating coherent trajectories from high-level instructions remains challenging especially for long-range composition tasks requiring multiple sequential skills. We propose SkillDiffuser an end-to-end hierarchical planning framework integrating interpretable skill learning with conditional diffusion planning to address this problem. At the higher level the skill abstraction module learns discrete human-understandable skill representations from visual observations and language instructions. These learned skill embeddings are then used to condition the diffusion model to generate customized latent trajectories aligned with the skills. This allows generating diverse state trajectories that adhere to the learnable skills. By integrating skill learning with conditional trajectory generation SkillDiffuser produces coherent behavior following abstract instructions across diverse tasks. Experiments on multi-task robotic manipulation benchmarks like Meta-World and LOReL demonstrate state-of-the-art performance and human-interpretable skill representations from SkillDiffuser. More visualization results and information could be found on https://skilldiffuser.github.io/.

Poster
SingularTrajectory: Universal Trajectory Predictor Using Diffusion Model

Inhwan Bae, Young-Jae Park, Hae-Gon Jeon

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There are five types of trajectory prediction tasks: deterministic stochastic domain adaptation momentary observation and few-shot. These associated tasks are defined by various factors such as the length of input paths data split and pre-processing methods. Interestingly even though they commonly take sequential coordinates of observations as input and infer future paths in the same coordinates as output designing specialized architectures for each task is still necessary. For the other task generality issues can lead to sub-optimal performances. In this paper we propose SingularTrajectory a diffusion-based universal trajectory prediction framework to reduce the performance gap across the five tasks. The core of SingularTrajectory is to unify a variety of human dynamics representations on the associated tasks. To do this we first build a Singular space to project all types of motion patterns from each task into one embedding space. We next propose an adaptive anchor working in the Singular space. Unlike traditional fixed anchor methods that sometimes yield unacceptable paths our adaptive anchor enables correct anchors which are put into a wrong location based on a traversability map. Finally we adopt a diffusion-based predictor to further enhance the prototype paths using a cascaded denoising process. Our unified framework ensures the generality across various benchmark settings such as input modality and trajectory lengths. Extensive experiments on five public benchmarks demonstrate that SingularTrajectory substantially outperforms existing models highlighting its effectiveness in estimating general dynamics of human movements. Code is publicly available at https://github.com/inhwanbae/SingularTrajectory.

Poster
Single Mesh Diffusion Models with Field Latents for Texture Generation

Thomas W. Mitchel, Carlos Esteves, Ameesh Makadia

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We introduce a framework for intrinsic latent diffusion models operating directly on the surfaces of 3D shapes with the goal of synthesizing high-quality textures. Our approach is underpinned by two contributions: Field Latents a latent representation encoding textures as discrete vector fields on the mesh vertices and Field Latent Diffusion Models which learn to denoise a diffusion process in the learned latent space on the surface. We consider a single-textured-mesh paradigm where our models are trained to generate variations of a given texture on a mesh. We show the synthesized textures are of superior fidelity compared those from existing single-textured-mesh generative models. Our models can also be adapted for user-controlled editing tasks such as inpainting and label-guided generation. The efficacy of our approach is due in part to the equivariance of our proposed framework under isometries allowing our models to seamlessly reproduce details across locally similar regions and opening the door to a notion of generative texture transfer. Code and visualizations are available at https://single-mesh-diffusion.github.io/.

Poster
SinSR: Diffusion-Based Image Super-Resolution in a Single Step

Yufei Wang, Wenhan Yang, Xinyuan Chen, Yaohui Wang, Lanqing Guo, Lap-Pui Chau, Ziwei Liu, Yu Qiao, Alex C. Kot, Bihan Wen

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While super-resolution (SR) methods based on diffusion models exhibit promising results their practical application is hindered by the substantial number of required inference steps. Recent methods utilize the degraded images in the initial state thereby shortening the Markov chain. Nevertheless these solutions either rely on a precise formulation of the degradation process or still necessitate a relatively lengthy generation path (e.g. 15 iterations). To enhance inference speed we propose a simple yet effective method for achieving single-step SR generation named SinSR. Specifically we first derive a deterministic sampling process from the most recent state-of-the-art (SOTA) method for accelerating diffusion-based SR. This allows the mapping between the input random noise and the generated high-resolution image to be obtained in a reduced and acceptable number of inference steps during training. We show that this deterministic mapping can be distilled into a student model that performs SR within only one inference step. Additionally we propose a novel consistency-preserving loss to simultaneously leverage the ground-truth image during the distillation process ensuring that the performance of the student model is not solely bound by the feature manifold of the teacher model resulting in further performance improvement. Extensive experiments conducted on synthetic and real-world datasets demonstrate that the proposed method can achieve comparable or even superior performance compared to both previous SOTA methods and the teacher model in just one sampling step resulting in a remarkable up to x10 speedup for inference. Our code will be released at https://github.com/wyf0912/SinSR/.

Poster
SimAC: A Simple Anti-Customization Method for Protecting Face Privacy against Text-to-Image Synthesis of Diffusion Models

Feifei Wang, Zhentao Tan, Tianyi Wei, Yue Wu, Qidong Huang

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Despite the success of diffusion-based customization methods on visual content creation increasing concerns have been raised about such techniques from both privacy and political perspectives. To tackle this issue several anti-customization methods have been proposed in very recent months predominantly grounded in adversarial attacks. Unfortunately most of these methods adopt straightforward designs such as end-to-end optimization with a focus on adversarially maximizing the original training loss thereby neglecting nuanced internal properties intrinsic to the diffusion model and even leading to ineffective optimization in some diffusion time steps. In this paper we strive to bridge this gap by undertaking a comprehensive exploration of these inherent properties to boost the performance of current anti-customization approaches. Two aspects of properties are investigated: 1) We examine the relationship between time step selection and the model's perception in the frequency domain of images and find that lower time steps can give much more contributions to adversarial noises. This inspires us to propose an adaptive greedy search for optimal time steps that seamlessly integrates with existing anti-customization methods. 2) We scrutinize the roles of features at different layers during denoising and devise a sophisticated feature-based optimization framework for anti-customization. Experiments on facial benchmarks demonstrate that our approach significantly increases identity disruption thereby protecting user privacy and copyright. Our code is available at: https://github.com/somuchtome/SimAC.

Poster
SiTH: Single-view Textured Human Reconstruction with Image-Conditioned Diffusion

Hsuan- I Ho, Jie Song, Otmar Hilliges

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A long-standing goal of 3D human reconstruction is to create lifelike and fully detailed 3D humans from single-view images. The main challenge lies in inferring unknown body shapes appearances and clothing details in areas not visible in the images. To address this we propose SiTH a novel pipeline that uniquely integrates an image-conditioned diffusion model into a 3D mesh reconstruction workflow. At the core of our method lies the decomposition of the challenging single-view reconstruction problem into generative hallucination and reconstruction subproblems. For the former we employ a powerful generative diffusion model to hallucinate unseen back-view appearance based on the input images. For the latter we leverage skinned body meshes as guidance to recover full-body texture meshes from the input and back-view images. SiTH requires as few as 500 3D human scans for training while maintaining its generality and robustness to diverse images. Extensive evaluations on two 3D human benchmarks including our newly created one highlighted our method's superior accuracy and perceptual quality in 3D textured human reconstruction. Our code and evaluation benchmark is available at https://ait.ethz.ch/sith.

Poster
Shadow Generation for Composite Image Using Diffusion Model

Qingyang Liu, Junqi You, Jianting Wang, Xinhao Tao, Bo Zhang, Li Niu

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In the realm of image composition generating realistic shadow for the inserted foreground remains a formidable challenge. Previous works have developed image-to-image translation models which are trained on paired training data. However they are struggling to generate shadows with accurate shapes and intensities hindered by data scarcity and inherent task complexity. In this paper we resort to foundation model with rich prior knowledge of natural shadow images. Specifically we first adapt ControlNet to our task and then propose intensity modulation modules to improve the shadow intensity. Moreover we extend the small-scale DESOBA dataset to DESOBAv2 using a novel data acquisition pipeline. Experimental results on both DESOBA and DESOBAv2 datasets as well as real composite images demonstrate the superior capability of our model for shadow generation task. The dataset code and model are released at https://github.com/bcmi/Object-Shadow-Generation-Dataset-DESOBAv2.

Poster