CVPR 2024 Papers

We are living in a world surrounded by diverse and "smart" devices with rich modalities of sensing ability. Conveniently capturing the interactions between us humans and these objects remains far-reaching. In this paper we present I'm-HOI a monocular scheme to faithfully capture the 3D motions of both the human and object in a novel setting: using a minimal amount of RGB camera and object-mounted Inertial Measurement Unit (IMU). It combines general motion inference and category-aware refinement. For the former we introduce a holistic human-object tracking method to fuse the IMU signals and the RGB stream and progressively recover the human motions and subsequently the companion object motions. For the latter we tailor a category-aware motion diffusion model which is conditioned on both the raw IMU observations and the results from the previous stage under over-parameterization representation. It significantly refines the initial results and generates vivid body hand and object motions. Moreover we contribute a large dataset with ground truth human and object motions dense RGB inputs and rich object-mounted IMU measurements. Extensive experiments demonstrate the effectiveness of I'm-HOI under a hybrid capture setting. Our dataset and code will be released to the community.

3D shape generation from text is a fundamental task in 3D representation learning. The text-shape pairs exhibit a hierarchical structure where a general text like "chair" covers all 3D shapes of the chair while more detailed prompts refer to more specific shapes. Furthermore both text and 3D shapes are inherently hierarchical structures. However existing Text2Shape methods such as SDFusion do not exploit that. In this work we propose HyperSDFusion a dual-branch diffusion model that generates 3D shapes from a given text. Since hyperbolic space is suitable for handling hierarchical data we propose to learn the hierarchical representations of text and 3D shapes in hyperbolic space. First we introduce a hyperbolic text-image encoder to learn the sequential and multi-modal hierarchical features of text in hyperbolic space. In addition we design a hyperbolic text-graph convolution module to learn the hierarchical features of text in hyperbolic space. In order to fully utilize these text features we introduce a dual-branch structure to embed text features in 3D feature space. At last to endow the generated 3D shapes with a hierarchical structure we devise a hyperbolic hierarchical loss. Our method is the first to explore the hyperbolic hierarchical representation for text-to-shape generation. Experimental results on the existing text-to-shape paired dataset Text2Shape achieved state-of-the-art results. We release our implementation under HyperSDFusion.github.io.

Personalization has emerged as a prominent aspect within the field of generative AI enabling the synthesis of individuals in diverse contexts and styles while retaining high-fidelity to their identities. However the process of personalization presents inherent challenges in terms of time and memory requirements. Fine-tuning each personalized model needs considerable GPU time investment and storing a personalized model per subject can be demanding in terms of storage capacity. To overcome these challenges we propose HyperDreamBooth - a hypernetwork capable of efficiently generating a small set of personalized weights from a single image of a person. By composing these weights into the diffusion model coupled with fast finetuning HyperDreamBooth can generate a person's face in various contexts and styles with high subject details while also preserving the model's crucial knowledge of diverse styles and semantic modifications. Our method achieves personalization on faces in roughly 20 seconds 25x faster than DreamBooth and 125x faster than Textual Inversion using as few as one reference image with the same quality and style diversity as DreamBooth. Also our method yields a model that is 10000x smaller than a normal DreamBooth model.

Generating a 3D human model from a single reference image is challenging because it requires inferring textures and geometries in invisible views while maintaining consistency with the reference image. Previous methods utilizing 3D generative models are limited by the availability of 3D training data. Optimization-based methods that lift text-to-image diffusion models to 3D generation often fail to preserve the texture details of the reference image resulting in inconsistent appearances in different views. In this paper we propose HumanRef a 3D human generation framework from a single-view input. To ensure the generated 3D model is photorealistic and consistent with the input image HumanRef introduces a novel method called reference-guided score distillation sampling (Ref-SDS) which effectively incorporates image guidance into the generation process. Furthermore we introduce region-aware attention to Ref-SDS ensuring accurate correspondence between different body regions. Experimental results demonstrate that HumanRef outperforms state-of-the-art methods in generating 3D clothed humans with fine geometry photorealistic textures and view-consistent appearances. Code and model are available at https://eckertzhang.github.io/HumanRef.github.io/.

Recent text-to-3D methods employing diffusion models have made significant advancements in 3D human generation. However these approaches face challenges due to the limitations of text-to-image diffusion models which lack an understanding of 3D structures. Consequently these methods struggle to achieve high-quality human generation resulting in smooth geometry and cartoon-like appearances. In this paper we propose HumanNorm a novel approach for high-quality and realistic 3D human generation. The main idea is to enhance the model's 2D perception of 3D geometry by learning a normal-adapted diffusion model and a normal-aligned diffusion model. The normal-adapted diffusion model can generate high-fidelity normal maps corresponding to user prompts with view-dependent and body-aware text. The normal-aligned diffusion model learns to generate color images aligned with the normal maps thereby transforming physical geometry details into realistic appearance. Leveraging the proposed normal diffusion model we devise a progressive geometry generation strategy and a multi-step Score Distillation Sampling (SDS) loss to enhance the performance of 3D human generation. Comprehensive experiments substantiate HumanNorm's ability to generate 3D humans with intricate geometry and realistic appearances. HumanNorm outperforms existing text-to-3D methods in both geometry and texture quality. The project page of HumanNorm is https://humannorm.github.io/.

Current subject-driven image generation methods encounter significant challenges in person-centric image generation. The reason is that they learn the semantic scene and person generation by fine-tuning a common pre-trained diffusion which involves an irreconcilable training imbalance. Precisely to generate realistic persons they need to sufficiently tune the pre-trained model which inevitably causes the model to forget the rich semantic scene prior and makes scene generation over-fit to the training data. Moreover even with sufficient fine-tuning these methods can still not generate high-fidelity persons since joint learning of the scene and person generation also lead to quality compromise. In this paper we propose Face-diffuser an effective collaborative generation pipeline to eliminate the above training imbalance and quality compromise. Specifically we first develop two specialized pre-trained diffusion models i.e. Text-driven Diffusion Model (TDM) and Subject-augmented Diffusion Model (SDM) for scene and person generation respectively. The sampling process is divided into three sequential stages i.e. semantic scene construction subject-scene fusion and subject enhancement. The first and last stages are performed by TDM and SDM respectively. The subject-scene fusion stage that is the collaboration achieved through a novel and highly effective mechanism Saliency-adaptive Noise Fusion (SNF). Specifically it is based on our key observation that there exists a robust link between classifier-free guidance responses and the saliency of generated images. In each time step SNF leverages the unique strengths of each model and allows for the spatial blending of predicted noises from both models automatically in a saliency-aware manner all of which can be seamlessly integrated into the DDIM sampling process. Extensive experiments confirm the impressive effectiveness and robustness of the Face-diffuser in generating high-fidelity person images depicting multiple unseen persons with varying contexts. Code is available at https://github.com/CodeGoat24/Face-diffuser.

Despite diffusion models having shown powerful abilities to generate photorealistic images generating videos that are realistic and diverse still remains in its infancy. One of the key reasons is that current methods intertwine spatial content and temporal dynamics together leading to a notably increased complexity of text-to-video generation (T2V). In this work we propose HiGen a diffusion model-based method that improves performance by decoupling the spatial and temporal factors of videos from two perspectives i.e. structure level and content level. At the structure level we decompose the T2V task into two steps including spatial reasoning and temporal reasoning using a unified denoiser. Specifically we generate spatially coherent priors using text during spatial reasoning and then generate temporally coherent motions from these priors during temporal reasoning. At the content level we extract two subtle cues from the content of the input video that can express motion and appearance changes respectively. These two cues then guide the model's training for generating videos enabling flexible content variations and enhancing temporal stability. Through the decoupled paradigm HiGen can effectively reduce the complexity of this task and generate realistic videos with semantics accuracy and motion stability. Extensive experiments demonstrate the superior performance of HiGen over the state-of-the-art T2V methods. We have released our source code and models.

Diffusion models have demonstrated remarkable performance in image and video synthesis. However scaling them to high-resolution inputs is challenging and requires restructuring the diffusion pipeline into multiple independent components limiting scalability and complicating downstream applications. In this work we study patch diffusion models (PDMs) -- a diffusion paradigm which models the distribution of patches rather than whole inputs keeping up to 0.7% of the original pixels. This makes it very efficient during training and unlocks end-to-end optimization on high-resolution videos. We improve PDMs in two principled ways. First to enforce consistency between patches we develop deep context fusion -- an architectural technique that propagates the context information from low-scale to high-scale patches in a hierarchical manner. Second to accelerate training and inference we propose adaptive computation which allocates more network capacity and computation towards coarse image details. The resulting model sets a new state-of-the-art FVD score of 66.32 and Inception Score of 87.68 in class-conditional video generation on UCF-101 256x256 surpassing recent methods by more than 100%. Then we show that it can be rapidly fine-tuned from a base 36x64 low-resolution generator for high-resolution 64x288x512 text-to-video synthesis. To the best of our knowledge our model is the first diffusion-based architecture which is trained on such high resolutions entirely end-to-end. Project webpage: https://snap-research.github.io/hpdm.

Recent progress in single-image 3D generation highlights the importance of multi-view coherency leveraging 3D priors from large-scale diffusion models pretrained on Internet-scale images. However the aspect of novel-view diversity remains underexplored within the research landscape due to the ambiguity in converting a 2D image into 3D content where numerous potential shapes can emerge. Here we aim to address this research gap by simultaneously addressing both consistency and diversity. Yet striking a balance between these two aspects poses a considerable challenge due to their inherent trade-offs. This work introduces HarmonyView a simple yet effective diffusion sampling technique adept at decomposing two intricate aspects in single-image 3D generation: consistency and diversity. This approach paves the way for a more nuanced exploration of the two critical dimensions within the sampling process. Moreover we propose a new evaluation metric based on CLIP image and text encoders to comprehensively assess the diversity of the generated views which closely aligns with human evaluators' judgments. In experiments HarmonyView achieves a harmonious balance demonstrating a win-win scenario in both consistency and diversity.

Extracting keypoint locations from input hand frames known as 3D hand pose estimation is a critical task in various human-computer interaction applications. Essentially the 3D hand pose estimation can be regarded as a 3D point subset generative problem conditioned on input frames. Thanks to the recent significant progress on diffusion-based generative models hand pose estimation can also benefit from the diffusion model to estimate keypoint locations with high quality. However directly deploying the existing diffusion models to solve hand pose estimation is non-trivial since they cannot achieve the complex permutation mapping and precise localization. Based on this motivation this paper proposes HandDiff a diffusion-based hand pose estimation model that iteratively denoises accurate hand pose conditioned on hand-shaped image-point clouds. In order to recover keypoint permutation and accurate location we further introduce joint-wise condition and local detail condition. Experimental results demonstrate that the proposed HandDiff significantly outperforms the existing approaches on four challenging hand pose benchmark datasets. Codes and pre-trained models are publicly available at https://github.com/cwc1260/HandDiff.

Reconstructing 3D hand mesh robustly from a single image is very challenging due to the lack of diversity in existing real-world datasets. While data synthesis helps relieve the issue the syn-to-real gap still hinders its usage. In this work we present HandBooster a new approach to uplift the data diversity and boost the 3D hand-mesh reconstruction performance by training a conditional generative space on hand-object interactions and purposely sampling the space to synthesize effective data samples. First we construct versatile content-aware conditions to guide a diffusion model to produce realistic images with diverse hand appearances poses views and backgrounds; favorably accurate 3D annotations are obtained for free. Then we design a novel condition creator based on our similarity-aware distribution sampling strategies to deliberately find novel and realistic interaction poses that are distinctive from the training set. Equipped with our method several baselines can be significantly improved beyond the SOTA on the HO3D and DexYCB benchmarks. Our code will be released on https://github.com/hxwork/HandBooster_Pytorch.

Text-to-image generative models can generate high-quality humans but realism is lost when generating hands. Common artifacts include irregular hand poses shapes incorrect numbers of fingers and physically implausible finger orientations. To generate images with realistic hands we propose a novel diffusion-based architecture called HanDiffuser that achieves realism by injecting hand embeddings in the generative process. HanDiffuser consists of two components: a Text-to-Hand-Params diffusion model to generate SMPL-Body and MANO-Hand parameters from input text prompts and a Text-Guided Hand-Params-to-Image diffusion model to synthesize images by conditioning on the prompts and hand parameters generated by the previous component. We incorporate multiple aspects of hand representation including 3D shapes and joint-level finger positions orientations and articulations for robust learning and reliable performance during inference. We conduct extensive quantitative and qualitative experiments and perform user studies to demonstrate the efficacy of our method in generating images with high-quality hands.

3D hand-object interaction data is scarce due to the hardware constraints in scaling up the data collection process. In this paper we propose HOIDiffusion for generating realistic and diverse 3D hand-object interaction data. Our model is a conditional diffusion model that takes both the 3D hand-object geometric structure and text description as inputs for image synthesis. This offers a more controllable and realistic synthesis as we can specify the structure and style inputs in a disentangled manner. HOIDiffusion is trained by leveraging a diffusion model pre-trained on large-scale natural images and a few 3D human demonstrations. Beyond controllable image synthesis we adopt the generated 3D data for learning 6D object pose estimation and show its effectiveness in improving perception systems. Project page: https://mq-zhang1.github.io/HOIDiffusion.

To date the quest to rapidly and effectively produce human-object interaction (HOI) animations directly from textual descriptions stands at the forefront of computer vision research. The underlying challenge demands both a discriminating interpretation of language and a comprehensive physics-centric model supporting real-world dynamics. To ameliorate this paper advocates HOIAnimator a novel and interactive diffusion model with perception ability and also ingeniously crafted to revolutionize the animation of complex interactions from linguistic narratives. The effectiveness of our model is anchored in two ground-breaking innovations: (1) Our Perceptive Diffusion Models (PDM) brings together two types of models: one focused on human movements and the other on objects. This combination allows for animations where humans and objects move in concert with each other making the overall motion more realistic. Additionally we propose a Perceptive Message Passing (PMP) mechanism to enhance the communication bridging the two models ensuring that the animations are smooth and unified; (2) We devise an Interaction Contact Field (ICF) a sophisticated model that implicitly captures the essence of HOIs. Beyond mere predictive contact points the ICF assesses the proximity of human and object to their respective environment informed by a probabilistic distribution of interactions learned throughout the denoising phase. Our comprehensive evaluation showcases HOIanimator's superior ability to produce dynamic context-aware animations that surpass existing benchmarks in text-driven animation synthesis.

Incorporating human feedback has been shown to be crucial to align text generated by large language models to human preferences. We hypothesize that state-of-the-art instructional image editing models where outputs are generated based on an input image and an editing instruction could similarly benefit from human feedback as their outputs may not adhere to the correct instructions and preferences of users. In this paper we present a novel framework to harness human feedback for instructional visual editing (HIVE). Specifically we collect human feedback on the edited images and learn a reward function to capture the underlying user preferences. We then introduce scalable diffusion model fine-tuning methods that can incorporate human preferences based on the estimated reward. Besides to mitigate the bias brought by the limitation of data we contribute a new 1.1M training dataset a 3.6K reward dataset for rewards learning and a 1K evaluation dataset to boost the performance of instructional image editing. We conduct extensive empirical experiments quantitatively and qualitatively showing that HIVE is favored over previous state-of-the-art instructional image editing approaches by a large margin.

Hyperspectral image (HSI) restoration aims at recovering clean images from degraded observations and plays a vital role in downstream tasks. Existing model-based methods have limitations in accurately modeling the complex image characteristics with handcraft priors and deep learning-based methods suffer from poor generalization ability. To alleviate these issues this paper proposes an unsupervised HSI restoration framework with pre-trained diffusion model (HIR-Diff) which restores the clean HSIs from the product of two low-rank components i.e. the reduced image and the coefficient matrix. Specifically the reduced image which has a low spectral dimension lies in the image field and can be inferred from our improved diffusion model where a new guidance function with total variation (TV) prior is designed to ensure that the reduced image can be well sampled. The coefficient matrix can be effectively pre-estimated based on singular value decomposition (SVD) and rank-revealing QR (RRQR) factorization. Furthermore a novel exponential noise schedule is proposed to accelerate the restoration process (about 5xacceleration for denoising) with little performance decrease. Extensive experimental results validate the superiority of our method in both performance and speed on a variety of HSI restoration tasks including HSI denoising noisy HSI super-resolution and noisy HSI inpainting. The code is available at https://github.com/LiPang/HIRDiff.

Recent years have witnessed a trend of the deep integration of the generation and reconstruction paradigms. In this paper we extend the ability of controllable generative models for a more comprehensive hand mesh recovery task: direct hand mesh generation inpainting reconstruction and fitting in a single framework which we name as Holistic Hand Mesh Recovery (HHMR). Our key observation is that different kinds of hand mesh recovery tasks can be achieved by a single generative model with strong multimodal controllability and in such a framework realizing different tasks only requires giving different signals as conditions. To achieve this goal we propose an all-in-one diffusion framework based on graph convolution and attention mechanisms for holistic hand mesh recovery. In order to achieve strong control generation capability while ensuring the decoupling of multimodal control signals we map different modalities to a share feature space and apply cross-scale random masking in both modality and feature levels. In this way the correlation between different modalities can be fully exploited during the learning of hand priors. Furthermore we propose Condition-aligned Gradient Guidance to enhance the alignment of the generated model with the control signals which significantly improves the accuracy of the hand mesh reconstruction and fitting. Experiments show that our novel framework can realize multiple hand mesh recovery tasks simultaneously and outperform the existing methods in different tasks which provides more possibilities for subsequent downstream applications including gesture recognition pose generation mesh editing and so on.

In this paper we present a method to reconstruct the world and multiple dynamic humans in 3D from a monocular video input. As a key idea we represent both the world and multiple humans via the recently emerging 3D Gaussian Splatting (3D-GS) representation enabling to conveniently and efficiently compose and render them together. In particular we address the scenarios with severely limited and sparse observations in 3D human reconstruction a common challenge encountered in the real world. To tackle this challenge we introduce a novel approach to optimize the 3D-GS representation in a canonical space by fusing the sparse cues in the common space where we leverage a pre-trained 2D diffusion model to synthesize unseen views while keeping the consistency with the observed 2D appearances. We demonstrate our method can reconstruct high-quality animatable 3D humans in various challenging examples in the presence of occlusion image crops few-shot and extremely sparse observations. After reconstruction our method is capable of not only rendering the scene in any novel views at arbitrary time instances but also editing the 3D scene by removing individual humans or applying different motions for each human. Through various experiments we demonstrate the quality and efficiency of our methods over alternative existing approaches.

Driven by the scalable diffusion models trained on large-scale datasets text-to-image synthesis methods have shown compelling results. However these models still fail to precisely follow the text prompt involving multiple objects attributes or spatial compositions. In this paper we reveal the potential causes of the diffusion model's cross-attention and self-attention layers. We propose two novel losses to refocus attention maps according to a given spatial layout during sampling. Creating the layouts manually requires additional effort and can be tedious. Therefore we explore using large language models (LLM) to produce these layouts for our method. We conduct extensive experiments on the DrawBench HRS and TIFA benchmarks to evaluate our proposed method. We show that our proposed attention refocusing effectively improves the controllability of existing approaches.